text:
# Untitled Document

Second-order linear differential equations with two irregular singular

points of rank three: the characteristic exponent

Wolfgang Bühring

Physikalisches Institut, Universität Heidelberg, Philosophenweg 12,

D-69120 Heidelberg, Germany

Abstract

For a second-order linear differential equation with two irregular singular points of rank three, multiple Laplace-type contour integral solutions are considered. An explicit formula in terms of the Stokes multipliers is derived for the characteristic exponent of the multiplicative solutions. The Stokes multipliers are represented by converging series with terms for which limit formulas as well as more detailed asymptotic expansions are available. Here certain new, recursively known coefficients enter , which are closely related to but different from the coefficients of the formal solutions at one of the irregular singular points of the differential equation. The coefficients of the formal solutions then appear as finite sums over subsets of the new coefficients. As a by-product, the leading exponential terms of the asymptotic behaviour of the late coefficients of the formal solutions are given, and this is a concrete example of the structural results obtained by Immink in a more general setting. The formulas displayed in this paper are not of merely theoretical interest, but they also are complete in the sense that they could be (and have been) implemented for computing accurate numerical values of the characteristic exponent, although the computational load is not small and increases with the rank of the singular point under consideration.

AMS classification: 34A20; 34A25; 34A30

Keywords: Irregular singular point; Characteristic exponent; Stokes multiplier

1. Introduction

Let us consider the differential equation

$$z^2f^{\prime \prime }+zf^{}[\underset{m=1}{\overset{6}{}}D_mz^m+L^2+\underset{m=1}{\overset{6}{}}B_mz^m]f(z)=0$$
$`(1.1)`$
with the thirteen parameters $`D_1,\mathrm{},D_6`$, $`L`$, $`B_1,\mathrm{},B_6`$, which for simplicity of presentation are assumed to be real. This differential equation has two irregular singular points, each of rank 3 (if $`D_60`$ and $`B_60`$, respectively), at the origin and at infinity. Without loss of generality, one of the parameters except $`L`$ could be set equal to $`1`$ . We assume that $`B_6`$ is positive, again for simplicity of presentation, and that $`L`$ is not negative.

At infinity, there are formal power series solutions

$$f_\mathrm{}1^{\mathrm{asy}}(z)=\mathrm{exp}(P(z))z^{\tau (1)}\underset{n=0}{\overset{\mathrm{}}{}}a_n(1)z^n,$$
$`(1.2)`$
$$f_\mathrm{}2^{\mathrm{asy}}(z)=\mathrm{exp}(P(z))z^{\tau (1)}\underset{n=0}{\overset{\mathrm{}}{}}a_n(1)z^n,$$
$`(1.3)`$
where the various quantities are determined by

$$P(z)=p_3z^3+p_2z^2+p_1z,$$
$`(1.4)`$
$$p_3=\frac{1}{3}\sqrt{B_6},p_2=\frac{1}{12}B_5/p_3,p_1=\frac{1}{6}(B_44p_2^2)/p_3,$$
$`(1.5)`$
$$\tau (\kappa )=\frac{3}{2}\frac{1}{6}\frac{B_3}{\kappa p_3}+\frac{2}{3}\frac{p_1p_2}{\kappa p_3},$$
$`(1.6)`$
$$\kappa \{1,1\}$$
$`(1.7)`$
and where the coefficients

$$a_n(\kappa )=a_n$$
$`(1.8)`$
are known recursively by

$$a_0=1$$
$`(1.9)`$
and

$$\begin{array}{cc}& 6\kappa p_3na_n=[4\kappa p_2(\tau (\kappa )+n2)+p_1^2B_2]a_{n1}\hfill \\ & +[2\kappa p_1(\tau (\kappa )+n\frac{5}{2})B_1]a_{n2}+(\tau (\kappa )+n3L)(\tau (\kappa )+n3+L)a_{n3}\hfill \\ & \underset{m=1}{\overset{6}{}}D_ma_{nm3},\hfill \end{array}$$
$`(1.10)`$
$$n=1,2,\mathrm{},(a_8=a_7=\mathrm{}=a_1=0).$$

The formal solutions are asymptotic expansions as $`z\mathrm{}`$ in appropriate sectors of the complex plane.

At the origin, there are analogical formal solutions, which may be obtained from those at infinity by the simultaneous replacements $`z1/z,B_mD_m,m=1,\mathrm{},6`$. This symmetry is our main reason for choosing just (1.1) as the standard form for a differential equation with two irregular singular points.

In the ring-shaped region $`0<|z|<\mathrm{}`$ we have (convergent) Floquet solutions $`f^{(\omega )}(z)`$ and $`f^{(\omega )}(z)`$, where

$$f^{(\omega )}(z)=z^\omega \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}c_n^{(\omega )}z^n,$$
$`(1.11)`$
which are linearly independent if $`2\omega `$ is not equal to an integer. Here the coefficients $`c_n^{(\omega )}`$ obey the recurrence relation

$$\underset{m=1}{\overset{6}{}}B_mc_{nm}^{(\omega )}+(L+\omega +n)(L+\omega +n)c_n^{(\omega )}\underset{m=1}{\overset{6}{}}D_mc_{n+m}^{(\omega )}=0,$$
$`(1.12)`$
and we want to normalize them by choosing

$$c_0^{(\omega )}=1.$$
$`(1.13)`$
The requirement that the power series converge determines, modulo 1, the possible values of the characteristic exponent (or circuit exponent or Floquet exponent) $`\omega `$. The problem to compute the characteristic exponent appears also and is best known in the context of Hill’s differential equation . There are methods to compute the characteristic exponent numerically, which require, for $`\omega =0`$, the evaluation of the infinite determinant  associated with (1.12), or numerical integration of the differential equation along a suitable contour , or numerical solution of an eigen-value problem .

This paper developes an entirely different method for evaluating the characteristic exponent. We obtain an explicit formula in terms of quantities which are essentially the Stokes multipliers, and these are given explicitly as convergent series, the terms of which are represented by asymptotic expansions. Here certain recursively known coefficients enter which are closely related to but different from the coefficients of the formal solutions (1.2)-(1.3).

Although some of the other authors concerned with irregular singular points of rank larger than one \[1, 2, 4-6, 8-10, 12, 13, 15, 18, 19\] consider the general case of arbitrary rank , we here prefer to restrict our attention to rank three. This is already general enough to give an impression of what can be expected in the case of even higher rank. On the other hand, it is still simple enough so that we can, for the relevant quantities, obtain explicit expressions which are not only theoretically interesting but can also be implemented (and have been implemented) for numerical evaluation. This work may be viewed as an attempt to extend , which was useful for rank one, to the much more complicated case of higher rank.

2. Laplace contour integral solutions

We try to apply the classical method of multiple Laplace contour integral solutions  and write for a solution of (1.1)

$$f(z)=z^\lambda (2\pi \text{i})^3_{C_{t_2}}_{C_{t_1}}_{C_s}\mathrm{exp}(z^3s+z^2t_1+zt_2)v(s,t_1,t_2)\text{d}s\text{d}t_1\text{d}t_2,$$
$`(2.1)`$
where a power factor with a still arbitrary parameter $`\lambda `$ has been included in view of later benefits. To derive the appropriate weight function $`v(s,t_1,t_2)`$ is a somewhat lengthy but not principally difficult procedure. We therefore give two lemmata stating the results of this procedure and postpone the proofs to a later section.

Lemma 1. The weight function $`v(s,t_1,t_2)`$ has to be a solution of the partial differential equation

$$\begin{array}{cc}& 9(s^2s_0^2)(^4v/s^4)+12(st_1s_0t_{10})(^4v/s^3t_1)+6(st_2s_0t_{20})(^4v/s^3t_2)\hfill \\ & +4(t_1^2t_{10}^2)(^4v/s^2t_1^2)+(t_2^2t_{20}^2)(^4v/s^2t_2^2)+4(t_1t_2t_{10}t_{20})(^4v/s^2t_1t_2)\hfill \\ & +([816\lambda ]s+[B_34t_{10}t_{20}])(^3v/s^3)+([524\lambda ]t_1+[B_2t_{20}^2])(^3v/s^2t_1)\hfill \\ & +([252\lambda ]t_2+B_1)(^3v/s^2t_2)+([\lambda 12]^2L^2)(^2v/s^2)D_1(^2v/st_1)\hfill \\ & D_2(^2v/st_2)+D_3(v/s)+D_4(v/t_1)+D_5(v/t_2)D_6v=0\hfill \end{array}$$
$`(2.2)`$
with two finite singular points at $`(s,t_1,t_2)=(\kappa s_0,\kappa t_{10},\kappa t_{20})`$ where

$$t_{20}=p_1,t_{10}=p_2,s_0=p_3,$$
$`(2.3)`$
and the contours for each variable have to satisfy the condition that a certain lengthy expression, bilinear in $`K`$ and $`v`$ or their partial derivatives, have the same value at the start and the end of the contour.

Lemma 2. For $`\kappa \{1,1\}`$, there are appropriate weight functions $`v=V(\kappa ;s,t_1,t_2)`$ which at the singular point $`(s,t_1,t_2)=(\kappa s_0,\kappa t_{10},\kappa t_{20})`$ have the power series expansion

$$\begin{array}{cc}& V(\kappa ;s,t_1,t_2)\hfill \\ & =\underset{m=0}{\overset{\mathrm{}}{}}\underset{n_1=0}{\overset{\mathrm{}}{}}\underset{n_2=0}{\overset{\mathrm{}}{}}A(\kappa ;m,n_1,n_2)(s\kappa s_0)^{\mu (\kappa )+m}(t_1\kappa t_{10})^{\nu _1n_1}(t_2\kappa t_{20})^{\nu _2n_2},\hfill \end{array}$$
$`(2.4)`$
where the $`A`$-coefficients satisfy a certain recurrence relation and the exponents are related by

$$3\mu (\kappa )2\nu _1\nu _2=\lambda +\tau (\kappa )6.$$
$`(2.5)`$
An appropriate set of coefficients is

$$A(\kappa ;m,n_1,n_2)=\mathrm{\Gamma }(\mu (\kappa )m)\mathrm{\Gamma }(\nu _1+n_1)\mathrm{\Gamma }(\nu _2+n_2)b(\kappa ;m,n_1,n_2),$$
$`(2.6)`$
where the new coefficients $`b(\kappa ;m,n_1,n_2)`$ are given by the recurrence relation

$$\begin{array}{cc}& 6\kappa s_0(3m2n_1n_2)b(\kappa ;m,n_1,n_2)\hfill \\ & =[(3m2n_1n_2+\tau (\kappa )3)^2L^2]b(\kappa ;m1,n_1,n_2)\hfill \\ & +[4\kappa t_{10}(3m2n_1n_2+\tau (\kappa )2)+t_{20}^2B_2]b(\kappa ;m1,n_11,n_2)\hfill \\ & +[2\kappa t_{20}(3m2n_1n_2+\tau (\kappa )\frac{5}{2})B_1]b(\kappa ;m1,n_1,n_21)\hfill \\ & D_1b(\kappa ;m2,n_11,n_2)D_2b(\kappa ;m2,n_1,n_21)D_3b(\kappa ;m2,n_1,n_2)\hfill \\ & D_4b(\kappa ;m3,n_11,n_2)D_5b(\kappa ;m3,n_1,n_21)D_6b(\kappa ;m3,n_1,n_2)\hfill \end{array}$$
$`(2.7)`$
with the initial conditions

$$b(\kappa ;0,0,0)=1,$$

$$b(\kappa ;m,n_1,n_2)=0\text{if}3m2n_1n_2=0\text{for}(m,n_1,n_2)(0,0,0),$$
$`(2.8)`$
$$(b(\kappa ;m,n_1,n_2)=0\text{if}m<0\text{or}n_1<0\text{or}n_2<0)$$

and with

$$b(\kappa ;0,n_1,n_2)=0\text{for}(n_1,n_2)(0,0),$$
$`(2.9)`$
$$b(\kappa ;m,n_1,n_2)=0\text{for}n_1+n_2>m$$
$`(2.10)`$
as a consequence. In addition, there are weight functions $`U(\kappa ;s,t_1,t_2)`$ which are analytic in $`s`$ at the respective singular point, corresponding to $`\mu (\kappa )=0,1,2`$ (without any further relation such as (2.5)).

In order to avoid unnecessary complications, we want to assume that the non-trivial exponent $`\mu `$ according to (2.5) is not equal to an integer. This can always be guarantied by a suitable choice of the still disposable parameter $`\lambda `$.

Since the exponents $`\mu (\kappa )`$, $`\nu _1`$, $`\nu _2`$ are restricted only by (2.5) but otherwise arbitrary, there are other solutions of the partial differential equation (2.2 ) relevant as weight functions in our contour integrals (2.1). We may assume that $`\nu _1`$ and $`\nu _2`$ are positive integers, preferentially

$$\nu _1=\nu _2=1,$$
$`(2.11)`$
but for the time being we want to keep $`\nu _1`$ and $`\nu _2`$ in the formulas. If $`\nu _1`$ and $`\nu _2`$ are increased by any positive integers $`q_1`$ and $`q_2`$, respectively, and $`\mu (\kappa )`$ simultaneously is decreased by $`(2/3)q_1+(1/3)q_2`$, then (2.5) is still satisfied. We therefore have to consider, for $`q_1,q_2=0,1,2,\mathrm{}`$, the set of solutions

$$\begin{array}{cc}& V(\kappa ;q_1,q_2)=V(\kappa ;q_1,q_2;s,t_1,t_2)\hfill \\ & =\underset{m=0}{\overset{\mathrm{}}{}}\underset{n_1=0}{\overset{m}{}}\underset{n_2=0}{\overset{m}{}}\mathrm{\Gamma }(\nu _1+q_1+n_1)\mathrm{\Gamma }(\nu _2+q_2+n_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2m)\hfill \\ & \times b(\kappa ;m,n_1,n_2)(s\kappa s_0)^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2+m}(t_1\kappa t_{10})^{\nu _1q_1n_1}(t_2\kappa t_{20})^{\nu _2q_2n_2}.\hfill \end{array}$$
$`(2.12)`$
As indicated, we will use a short-hand notation suppressing the dependence on the variables $`s,t_1,t_2`$. We now have to choose appropriate contours for the integral representation (2.1). For each of the integrals over $`t_1`$ or $`t_2`$ a closed circle, traversed once in the positive sense, around the relevant singular point of the integrand is appropriate, since the pertinent exponent is an integer. For the $`s`$ -integral we need an infinite contour which starts somewhere at infinity where the exponential factor of the integral vanishes, surrounds one of the singular points in the positive sense, and returns (on a different sheet) to the starting point. Assuming that $`s_0`$ is real and positive, then, if

$$0<\mathrm{arg}(z)<\frac{1}{3}\pi $$
$`(2.13)`$
the starting- and end-point at infinity has the phase $`\pi /2`$. If we agree that

$$\mathrm{arg}(s\kappa s_0)=0$$
$`(2.14)`$
when $`s`$ is positive and sufficiently large, the integral of a single term of the infinite series (2.4) can be evaluated:

$$\begin{array}{cc}& z^\lambda (2\pi \text{i})^3\stackrel{(\kappa t_{20}+)}{}\stackrel{(\kappa t_{10}+)}{}\underset{\kappa s_0+\mathrm{i}\mathrm{}}{\overset{(\kappa s_0+)}{}}\mathrm{exp}(z^3s+z^2t_1+zt_2)\hfill \\ & \times (s\kappa s_0)^{\mu (\kappa )+m}(t_1\kappa t_{10})^{\nu _1n_1}(t_2\kappa t_{20})^{\nu _2n_2}\text{d}s\text{d}t_1\text{d}t_2\hfill \\ & =\mathrm{exp}(\kappa s_0z^3+\kappa t_{10}z^2+\kappa t_{20}z)\frac{z^{\tau (\kappa )3m+2n_1+n_2}}{\mathrm{\Gamma }(\mu (\kappa )m)\mathrm{\Gamma }(\nu _1+n_1)\mathrm{\Gamma }(\nu _2+n_2)},\hfill \end{array}$$
$`(2.15)`$
where, in the power of $`z`$ on the right-hand side, already use has been made of (2.5). As a consequence, the integral of $`V(\kappa ;q_1,q_2)`$ yields, if the series is integrated term by term, one or the other of the formal solutions (1.2)-(1.3):

$$\begin{array}{cc}\hfill f_\mathrm{}j(z):=& z^\lambda (2\pi \text{i})^3\stackrel{(\kappa t_{20}+)}{}\stackrel{(\kappa t_{10}+)}{}\underset{\kappa s_0+\mathrm{i}\mathrm{}}{\overset{(\kappa s_0+)}{}}\mathrm{exp}(z^3s+z^2t_1+zt_2)\hfill \\ & \times V(\kappa ;q_1,q_2;s,t_1,t_2)\text{d}s\text{d}t_1\text{d}t_2f_\mathrm{}j^{\mathrm{asy}}(z)\hfill \end{array}$$
$`(2.16)`$
in the sector $`0<\mathrm{arg}(z)<\frac{1}{3}\pi `$, where $`j=1`$ if $`\kappa =1`$ or $`j=2`$ if $`\kappa =1`$ . Each of the solutions defined by the integral representation (2.16 ) has one of the formal solutions (1.2) or (1.3) as its asymptotic expansion as $`z\mathrm{}`$ in the indicated sector. It follows by rotation of the contour that the asymptotic expansions are theoretically valid in the larger sectors $`\frac{1}{6}\pi <\mathrm{arg}(z)<\frac{5}{6}\pi `$ for $`j=1`$ or $`\frac{1}{2}\pi <\mathrm{arg}(z)<\frac{1}{2}\pi `$ for $`j=2`$, respectively.

Looking at (2.15) we may see that all the terms for which $`3m2n_1n_2`$ is the same yield the same power of $`z`$. For the coefficients of the asymptotic expansions (1.2) or (1.3) we then have the representation

$$a_n(\kappa )=\underset{(m,n_1,n_2)I_n}{}b(\kappa ;m,n_1,n_2),$$
$`(2.17)`$
where

$$I_n=\{(m,n_1,n_2):3m2n_1n_2=n\}\text{N}_0\times \text{N}_0\times \text{N}_0,n=0,1,\mathrm{}.$$
$`(2.18)`$
Because of the properties of the $`b(m,n_1,n_2)`$, the sum in (2.18) is finite for each finite $`n`$, in particular we have

$$a_0(\kappa )=b_0(\kappa ;0,0,0)=1.$$
$`(2.19)`$
3. Analytic continuation of the integrand

Below we have to consider the integral representation (2.1) with $`t_1`$ and $`t_2`$-contours which are simple closed curves surrounding in the positive sense both the finite singular points $`t_{10}`$ and $`t_{10}`$ or $`t_{20}`$ and $`t_{20}`$, respectively, and with an $`s`$-contour which starts at or near $`s_0+\text{i}\mathrm{}`$, surrounds both the finite singular points $`s_0`$ and $`s_0`$ once in the positive sense and ends at $`s_0+\text{i}\mathrm{}`$. We therefore need the analytic continuation of the integrand between the two singular points along this contour. With appropriate power factors $`\mathrm{\Phi }`$ included for later convenience, the continuation formula reads

$$\begin{array}{cc}& \mathrm{\Phi }(\kappa ;r_1,r_2)V(\kappa ;r_1,r_2)\hfill \\ & =\underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Phi }(\kappa ;q_1,q_2)V(\kappa ;q_1,q_2)+U(\kappa ;r_1,r_2),\hfill \end{array}$$
$`(3.1)`$
where

$$\mathrm{\Phi }(\kappa ;r_1,r_2)=(2\kappa s_0)^{\mu (\kappa ){\scriptscriptstyle \frac{2}{3}}r_1{\scriptscriptstyle \frac{1}{3}}r_2}.$$
$`(3.2)`$
The effect of these power factors is that in (3.1) the total powers with non-integer exponents are powers of $`\frac{1}{2}[1s/(\kappa s_0)]`$ or of $`\frac{1}{2}[1+s/(\kappa s_0)]`$, respectively. We may agree that $`\mathrm{arg}(1s/s_0)=\mathrm{arg}(1+s/s_0)=0`$ when $`s`$ is on the real axis between $`s_0`$ and $`s_0`$. Also, as above, $`\mathrm{arg}(s)=0`$ when s is larger than $`s_0`$. Then, by analytic continuation along the contour under consideration, we have

$$\mathrm{\Phi }(1;r_1,r_2)=(2s_0)^{\mu (1){\scriptscriptstyle \frac{2}{3}}r_1{\scriptscriptstyle \frac{1}{3}}r_2}\mathrm{exp}(\text{i}\pi (\mu (1)+\frac{2}{3}r_1+\frac{1}{3}r_2),$$
$`(3.3)`$
$$\mathrm{\Phi }(1;r_1,r_2)=(2s_0)^{\mu (1){\scriptscriptstyle \frac{2}{3}}r_1{\scriptscriptstyle \frac{1}{3}}r_2}.$$
$`(3.4)`$
Let us rewrite (3.1) using an even more condensed notation, writing $`q`$ for $`(q_1,q_2)`$ in the parameter list of the various functions and writing a sum over $`q`$ in place of a double sum over $`q_1`$ and $`q_2`$, etc. The above continuation formula then reads

$$\mathrm{\Phi }(\kappa ;r)V(\kappa ;r)=\underset{q=r}{}E(\kappa ;r,q)\mathrm{\Phi }(\kappa ;q)V(\kappa ;q)+U(\kappa ;r),$$
$`(3.5)`$
and the second continuation formula

$$\begin{array}{cc}& U(\kappa ;r)=\mathrm{\Phi }(\kappa ,r)V(\kappa ,r)\hfill \\ & \underset{q=r}{}\underset{p=q}{}E(\kappa ;r,q)E(\kappa ;q,p)\mathrm{\Phi }(\kappa ;p)V(\kappa ;p)\underset{q=r}{}E(\kappa ;r,q)U(\kappa ;q)\hfill \end{array}$$
$`(3.6)`$
follows by the requirement of consistency of (3.5) for $`\kappa `$ replaced by $`\kappa `$.

4. Asymptotic expansions for the coefficients in the continuation formula

We now want to determine the $`E`$-coefficients in the continuation formulas (3.5) , (3.6) by means of the asymptotic method of Darboux  applied to the variable $`s`$ . The left-hand side of the continuation formula (3.5) is

$$\begin{array}{cc}& (\frac{1}{2}\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}r_1+{\scriptscriptstyle \frac{1}{3}}r_2}\underset{m=0}{\overset{\mathrm{}}{}}\underset{n_1=0}{\overset{m}{}}\underset{n_2=0}{\overset{m}{}}\mathrm{\Gamma }(\nu _1+r_1+n_1)\mathrm{\Gamma }(\nu _2+r_2+n_2)\hfill \\ & \times \mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2m)b(\kappa ;m,n_1,n_2)(2\kappa s_0)^m\hfill \\ & \times (\frac{1}{2}\frac{s}{2\kappa s_0})^m(t_1\kappa t_{10})^{\nu _1r_1n_1}(t_2\kappa t_{20})^{\nu _2r_2n_2}.\hfill \end{array}$$
$`(4.1)`$
The leading singular term, when $`(s,t_1,t_2)(\kappa s_0,\kappa t_{10},\kappa t_{20})`$, on the right-hand side is

$$\begin{array}{cc}& \underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \mathrm{\Gamma }(\nu _1+q_1)\mathrm{\Gamma }(\nu _2+q_2)b(\kappa ;0,0,0)\hfill \\ & \times (\frac{1}{2}+\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2}(t_1+\kappa t_{10})^{\nu _1q_1}(t_2+\kappa t_{20})^{\nu _2q_2}.\hfill \end{array}$$
$`(4.2)`$
By means of the binomial theorem in its hypergeometric-series-form

$$(1x)^\alpha =\underset{j=0}{\overset{\mathrm{}}{}}\frac{(\alpha )_j}{j!}x^j,$$
$`(4.3)`$
where

$$(\alpha )_j=\alpha (\alpha +1)\mathrm{}(\alpha +j1)=\mathrm{\Gamma }(\alpha +j)/\mathrm{\Gamma }(\alpha )$$

means the Pochhammer symbol, we may expand, if $`|t|`$ is sufficiently large,

$$(t+\kappa t_0)^{\nu q}=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(\nu +q)_j}{j!}(2\kappa t_0)^j(t\kappa t_0)^{\nu qj}$$
$`(4.4)`$
and, if $`|s/(\kappa s_0)|`$ is sufficiently small,

$$(\frac{1}{2}+\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2}=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)_m}{m!}(\frac{1}{2}\frac{s}{2\kappa s_0})^m.$$
$`(4.5)`$
Then the leading singular term on the right becomes

$$\begin{array}{cc}& \underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \mathrm{\Gamma }(\nu _1+q_1)\underset{j_1=0}{\overset{\mathrm{}}{}}\frac{(\nu _1+q_1)_{j_1}}{j_1!}(2\kappa t_{10})^{j_1}(t_1\kappa t_{10})^{\nu _1q_1j_1}\hfill \\ & \times \mathrm{\Gamma }(\nu _2+q_2)\underset{j_2=0}{\overset{\mathrm{}}{}}\frac{(\nu _2+q_2)_{j_2}}{j_2!}(2\kappa t_{20})^{j_2}(t_2\kappa t_{20})^{\nu _2q_2j_2}\hfill \\ & \times b(\kappa ;0,0,0)\underset{m=0}{\overset{\mathrm{}}{}}\frac{(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)_m}{m!}(\frac{1}{2}\frac{s}{2\kappa s_0})^m.\hfill \end{array}$$
$`(4.6)`$
The coefficients of this series should agree asymptotically, as $`m\mathrm{}`$, with those of the series on the left-hand side, where the power factor in front of the series (4.1), when $`s\kappa s_0`$, tends to $`1`$ and may be omitted. We therefore obtain

$$\begin{array}{cc}& \underset{n_1=0}{\overset{m}{}}\underset{n_2=0}{\overset{m}{}}\mathrm{\Gamma }(\nu _1+r_1+n_1)\mathrm{\Gamma }(\nu _2+r_2+n_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2m)\hfill \\ & \times b(\kappa ;m,n_1,n_2)(2\kappa s_0)^m(t_1\kappa t_{10})^{\nu _1r_1n_1}(t_2\kappa t_{20})^{\nu _2r_2n_2}\hfill \\ & \underset{q_1=r_1}{\overset{\mathrm{}}{}}\underset{q_2=r_2}{\overset{\mathrm{}}{}}E(\kappa ;r_1,r_2;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \mathrm{\Gamma }(\nu _1+q_1)\underset{j_1=0}{\overset{\mathrm{}}{}}\frac{(\nu _1+q_1)_{j_1}}{j_1!}(2\kappa t_{10})^{j_1}(t_1\kappa t_{10})^{\nu _1q_1j_1}\hfill \\ & \times \mathrm{\Gamma }(\nu _2+q_2)\underset{j_2=0}{\overset{\mathrm{}}{}}\frac{(\nu _2+q_2)_{j_2}}{j_2!}(2\kappa t_{20})^{j_2}(t_2\kappa t_{20})^{\nu _2q_2j_2}\hfill \\ & \times b(\kappa ;0,0,0)\frac{(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)_m}{m!},\hfill \end{array}$$
$`(4.7)`$
which holds asymptotically as $`m\mathrm{}`$. On both sides of this asymptotic equation a double power series in the same variables appears, so the coefficients of the corresponding terms must be equal. This yields, for each set $`r_1,r_2,n_1,n_2`$,

$$\begin{array}{cc}& \mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2m)(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{q_1=r_1}{\overset{r_1+n_1}{}}\underset{q_2=r_2}{\overset{r_2+n_2}{}}E(\kappa ;r_1,r_2;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{m!}\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=r_1+n_1q_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=r_2+n_2q_2}b(\kappa ;0,0,0)\hfill \end{array}$$
$`(4.8)`$
or, if we introduce new indices of summation and make use of the reflection formula of the gamma function,

$$\begin{array}{cc}& \pi [\mathrm{\Gamma }(1+m)]^1(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}E(\kappa ;r_1,r_2;r_1+p_1,r_2+p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}r_1\frac{1}{3}r_2\frac{2}{3}p_1\frac{1}{3}p_2+m)}{\mathrm{\Gamma }(1+m)}\hfill \\ & \times \frac{\mathrm{\Gamma }(1+\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2+m)}{\mathrm{\Gamma }(1+m)}\mathrm{sin}(\pi [\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2])\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=n_1p_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=n_2p_2}b(\kappa ;0,0,0).\hfill \end{array}$$
$`(4.9)`$
Now the left-hand side is independent of $`r_1`$ and $`r_2`$ and so is the product of the two ratios of gamma functions, when $`m\mathrm{}`$, on the right. Therefore $`E(\kappa ;r_1,r_2;r_1+p_1,r_2+p_2)\mathrm{sin}(\pi [\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2])`$ must be independent of $`r_1`$ and $`r_2`$ too. This proves

Lemma 3.

$$E(\kappa ;r_1,r_2;r_1+p_1,r_2+p_2)=\frac{\mathrm{sin}(\pi \mu (\kappa ))}{\mathrm{sin}(\pi [\mu (\kappa )+\frac{2}{3}r_1+\frac{1}{3}r_2])}E(\kappa ;0,0;p_1,p_2).$$
$`(4.10)`$
After Lemma 3 is available, we need to determine the $`E`$-coefficients for $`r_1=r_2=0`$ only, and it is advantageous to introduce the closely related coefficients

$$e(\kappa ;p_1,p_2)=\mathrm{sin}(\pi \mu (\kappa ))E(\kappa ;0,0;p_1,p_2),$$
$`(4.11)`$
which are independent of $`\lambda `$. Eq (4.9), with $`r_1=r_2=0`$ and the factor $`b(\kappa ;0,0,0)`$, which is equal to $`1`$, omitted, then becomes

$$\begin{array}{cc}& \frac{\pi }{\mathrm{\Gamma }(1+\mu (\kappa )+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+m)}{\mathrm{\Gamma }(1+m)}\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=n_1p_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=n_2p_2}.\hfill \end{array}$$
$`(4.12)`$
We then may solve (4.12) for the $`e`$-coefficient with $`p_1=n_1,p_2=n_2`$ and obtain the following asymptotic formula in terms of a $`b`$-coefficient and the earlier $`e`$-coefficients,

$$\begin{array}{cc}& e(\kappa ;n_1,n_2)\pi (2\kappa s_0)^m\hfill \\ & \times \frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}b(\kappa ;m,n_1,n_2)\hfill \\ & \underset{(p_1,p_2)(n_1,n_2)}{\underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}}e(\kappa ;p_1,p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+m)}{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1=n_1p_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2=n_2p_2}.\hfill \end{array}$$
$`(4.13)`$
By repeated application of this formula, all the $`e`$-coefficients on the right-hand side may be eliminated, and this yields the remarkable explicit limit formula

$$\begin{array}{cc}& e(\kappa ;n_1,n_2)=\pi \hfill \\ & \times \underset{m\mathrm{}}{lim}\frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times (2\kappa s_0)^m\underset{j_1=0}{\overset{n_1}{}}\underset{j_2=0}{\overset{n_2}{}}b(\kappa ;m,j_1,j_2)\frac{(2\kappa t_{10})^{n_1j_1}}{(n_1j_1)!}\frac{(2\kappa t_{20})^{n_2j_2}}{(n_2j_2)!}.\hfill \end{array}$$
$`(4.14)`$
The proof of this formula will be given below in Section 7.

Although approximate asymptotic equations or limit formulas such as (4.13) or (4.14) are interesting from a theoretical point of view, we finally need more, namely a detailed asymptotic expansion suitable for accurate numerical evaluation. This can be obtained, on the basis of Schäfke and Schmidt , essentially in the same way as above, apart from the following two refinements: The power factor in front of (4.1) can no longer be omitted, and we have to include a finite number of singular terms rather than the leading one, (4.2), alone. The result of this procedure, which will be derived in more detail below, is

Theorem 1. The $`E`$-coefficients, or $`e`$-coefficients according to (4.11), in the continuation formula (3.1) or (3.5) are

$$\begin{array}{cc}& e(\kappa ;n_1,n_2)=\pi \hfill \\ & \times [\frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{(q_1,q_2)(n_1,n_2)}{\underset{q_1=0}{\overset{n_1}{}}\underset{q_2=0}{\overset{n_2}{}}}e(\kappa ;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times [1+\underset{k=1}{\overset{K}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_k}{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2m)_k}H(\kappa ;k,l_1,l_2;q_1,q_2)+O(m^{K1})]\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1+q_1+l_1=n_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2+q_2+l_2=n_2}]\hfill \\ & \times \left[1+\underset{k=1}{\overset{K}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2)_k}{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2m)_k}H(\kappa ;k,0,0;n_1,n_2)+O(m^{K1})\right]^1,\hfill \end{array}$$
$`(4.15)`$
where

$$H(\kappa ;k,l_1,l_2;q_1,q_2)=\underset{j=l_1+l_2}{\overset{k}{}}\frac{(\mu (\kappa ))_{kj}}{(kj)!(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_j}(2\kappa s_0)^jb(\kappa ;j,l_1,l_2).$$
$`(4.16)`$
With a suitable choice of $`m`$ and $`K`$, Theorem 1 may be used to compute accurate values of the $`e`$-coefficients, beginning with

$$\begin{array}{cc}& e(\kappa ;0,0)=\pi \frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )+m)}(2\kappa s_0)^mb(\kappa ;m,0,0)\hfill \\ & \times \left[1+\underset{k=1}{\overset{K}{}}\frac{(1+\mu (\kappa ))_k}{(1+\mu (\kappa )m)_k}H(\kappa ;k,0,0;0,0)+O(m^{K1})\right]^1.\hfill \end{array}$$
$`(4.17)`$
While the $`e`$-coefficients are independent of $`\lambda `$, the approximate values computed for any finite $`m`$ do depend on it. Thus $`\lambda `$ here plays the role of a computational parameter which could be adjusted for optimal accuracy. The best choice from this point of view , in particular when $`L`$ is not small, is $`\lambda =L`$ or $`\lambda =L`$, according to the discussion in Section 7 below.

5. Multiplicative solutions and the Floquet exponent

We now want to construct a linear combination of the solutions at infinity,

$$f_p(z)=\alpha f_\mathrm{}1(z)+\beta f_\mathrm{}2(z),$$
$`(5.1)`$
which is a multiplicative solution such that, after analytic continuation along a sufficiently large circle around the origin traversed once in the negative sense, this solution remains the same apart from multiplication by a constant factor, that is

$$f_p(\mathrm{e}^{2\pi \mathrm{i}}z)=pf_p(z).$$
$`(5.2)`$
This solution is proportional to one of the Floquet solutions introduced above in the introduction.

We use the integral representation with contours which surround both the finite singular points as introduced above in Section 3. This integral is equal to the sum of the two integrals with contours, also considered above in Section 2, which surround only one of these singular points. If the contour is kept fixed, then the circle in the $`z`$-plane traversed once in the negative sense corresponds to a circle, around the origin and with radius greater than $`s_0`$, in the $`s`$-plane traversed three times in the positive sense, since then $`z^3s`$ in the exponential part of the integrand does not change.

Let us consider the integral representation with $`V(1;0)`$ in the integrand and with the above phase conventions. From (3.5) we have

$$\mathrm{\Phi }(1;0)V(1;0)=\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)V(1;q)+U(1;0).$$
$`(5.3)`$
The integral then yields a function

$$f^\mathrm{I}(z):=\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)f_\mathrm{}2(z),$$
$`(5.4)`$
where the two terms come from the two singular points, which contribute via the left- or right-hand side of (5.3), respectively. Let us now consider, in the $`s`$-plane, a simple closed loop, homotopic to the circle mentioned above, consisting of small circles around the singular points $`s_0`$ and $`s_0`$ and straight line segments along the real axis between these singular points. Let us start with the continuation formula (5.3) in a neighbourhood of the origin, where the left- as well as the right-hand side of (5.3) are valid, and see what happens when we follow the loop in the positive sense. We shall always refer to the above phase conventions and display any additional phases explicitly. Traversing the circular part around $`s_0`$ multiplies the function $`V(1;0)`$ by the phase factor $`\mathrm{exp}(2\pi \text{i}\mu (1))`$, so that we then have

$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}\mu (1))\mathrm{\Phi }(1;0)V(1;0)\hfill \\ & =\mathrm{exp}(2\pi \text{i}\mu (1))\{\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)V(1;q)+U(1;0)\},\hfill \end{array}$$
$`(5.5)`$
where the right-hand side is the analytic continuation of the left by means of (5.3). Next we have to traverse the circular part around $`s_0`$, which multiplies the function $`V(1;q)`$ by the phase factor $`\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])`$, where

$$\phi (q)=\frac{2}{3}q_1+\frac{1}{3}q_2,$$
$`(5.6)`$
so that we have on the straight line segment from $`s_0`$ to $`s_0`$ , after the loop has been traversed once,

$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}[\mu (1)+\mu (1)])\underset{q=0}{}E(1;0,q)\mathrm{exp}(2\pi \text{i}\phi (q))\mathrm{\Phi }(1;q)V(1;q)\hfill \\ & +\mathrm{exp}(2\pi \text{i}\mu (1))U(1;0)\hfill \end{array}$$

$$\begin{array}{cc}& =\mathrm{exp}(2\pi \text{i}\mu (1))\{\mathrm{\Phi }(1;0)V(1;0)+\underset{q=0}{}\{\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])1\}U(1;q)\hfill \\ & +\underset{q=0}{}\underset{p=q}{}E(1;0,q)\{\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])1\}E(1;q,p)\mathrm{\Phi }(1;p)V(1;p)\},\hfill \end{array}$$
$`(5.7)`$
where again the right-hand side is the analytic continuation of the left. It is evident that following the loop further leads to increasingly lengthier and more complicated formulas, not suitable for being fully displayed in this paper. We therefore want to stop here for a moment and consider the integral representation with the integrand obtained after the loop in the $`s`$-plane has been traversed only once (rather than three times, as finally needed). A representative example of the terms in (5.7) then is

$$\underset{q=0}{}\underset{p=q}{}E(1;0,q)\{\mathrm{exp}(2\pi \text{i}[\mu (1)+\phi (q)])1\}E(1;q,p)\mathrm{\Phi }(1;p)V(1;p).$$
$`(5.8)`$
The integral of $`V(1;p)`$ yields $`f_\mathrm{}1(z)`$, independent of $`p`$, so that the sum over $`p`$ can now be performed and , because of Lemma 3, the multiple sum reduces to a product of single sums. After integration we therefore obtain for (5.8)

$$\begin{array}{cc}& \{2\text{i}\mathrm{sin}(\pi \mu (1))\mathrm{exp}(\text{i}\pi \mu (1))\underset{q=0}{}E(1;0,q)(2s_0)^{\phi (q)}\mathrm{exp}(2\pi \text{i}\phi (q))\hfill \\ & \times \underset{\stackrel{~}{p}=0}{}E(1;0,\stackrel{~}{p})(2s_0)^{\phi (\stackrel{~}{p})}\mathrm{exp}(\text{i}\pi \phi (\stackrel{~}{p}))\}\mathrm{\Phi }(1;0)f_\mathrm{}1(z).\hfill \end{array}$$
$`(5.9)`$
It is now convenient to introduce the phase factor

$$\eta :=\mathrm{exp}(\frac{1}{3}\pi \text{i})=\frac{1}{2}(1+\text{i}\sqrt{3})$$
$`(5.10)`$
which satisfies

$$\eta ^6=1,1+\eta ^2+\eta ^4=0.$$
$`(5.11)`$
Also, for the sums occurring here and below we may introduce the ”Stokes multipliers”

$$\begin{array}{cc}& \sigma _n(1):=\mathrm{sin}(\pi \mu (1))\underset{q=0}{}E(1;0,q)(2s_0)^{\phi (q)}\mathrm{exp}(2n\pi \text{i}\phi (q)),\hfill \\ & \sigma _n(1):=\mathrm{sin}(\pi \mu (1))\underset{q=0}{}E(1;0,q)(2s_0)^{\phi (q)}\mathrm{exp}((2n+1)\pi \text{i}\phi (q)),\hfill \end{array}$$
$`(5.12)`$
which become

$$\begin{array}{cc}& \sigma _0(1)=S_0(1)+(2s_0)^{1/3}S_1(1)+(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _0(1)=S_0(1)+\eta (2s_0)^{1/3}S_1(1)+\eta ^2(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _1(1)=S_0(1)+\eta ^2(2s_0)^{1/3}S_1(1)+\eta ^4(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _1(1)=S_0(1)+\eta ^3(2s_0)^{1/3}S_1(1)+(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _2(1)=S_0(1)+\eta ^4(2s_0)^{1/3}S_1(1)+\eta ^2(2s_0)^{2/3}S_2(1),\hfill \\ & \sigma _2(1)=S_0(1)+\eta ^5(2s_0)^{1/3}S_1(1)+\eta ^4(2s_0)^{2/3}S_2(1)\hfill \end{array}$$
$`(5.13)`$
in terms of the real partial sums

$$\begin{array}{cc}\hfill S_0(\kappa ):=& \underset{l=0}{\overset{\mathrm{}}{}}(2\kappa s_0)^l\underset{(n_1,n_2)J_{3l}}{}e(\kappa ;n_1,n_2)=e(\kappa ;0,0)+\mathrm{},\hfill \\ \hfill S_1(\kappa ):=& \underset{l=0}{\overset{\mathrm{}}{}}(2\kappa s_0)^l\underset{(n_1,n_2)J_{3l+1}}{}e(\kappa ;n_1,n_2)=e(\kappa ;0,1)+\mathrm{},\hfill \\ \hfill S_2(\kappa ):=& \underset{l=0}{\overset{\mathrm{}}{}}(2\kappa s_0)^l\underset{(n_1,n_2)J_{3l+2}}{}e(\kappa ;n_1,n_2)=e(\kappa ;1,0)+e(\kappa ;0,2)+\mathrm{},\hfill \end{array}$$
$`(5.14)`$
where

$$J_l=\{(n_1,n_2):2n_1+n_2=l\}\text{N}_0\times \text{N}_0,l=0,1,\mathrm{}.$$
$`(5.15)`$
The integral of the representative term (5.9) then becomes

$$\frac{2\text{i}\mathrm{exp}(\text{i}\pi \mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _0(1)\mathrm{\Phi }(1;0)f_\mathrm{}1(z).$$
$`(5.16)`$
In total, we have

$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}\lambda /3)f^\mathrm{I}(\mathrm{e}^{2\pi \mathrm{i}/3}z)=\frac{\mathrm{exp}(2\pi \text{i}[\mu (1)+\mu (1)])}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\mathrm{\Phi }(1;0)f_\mathrm{}2(z)\hfill \\ & +\mathrm{exp}(2\pi \text{i}\mu (1)\{1+\frac{2\text{i}\mathrm{exp}(\pi \text{i}\mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _0(1)\}\mathrm{\Phi }(1;0)f_\mathrm{}1(z),\hfill \end{array}$$
$`(5.17)`$
where

$$f^\mathrm{I}(z)=\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+\frac{1}{\mathrm{sin}(\pi \mu (1))}\sigma _0(1)\mathrm{\Phi }(1;0)f_\mathrm{}2(z).$$
$`(5.18)`$
So far we have traversed the $`s`$-loop once and obtained on it the analytic continuation of the integrand, but we have to traverse it three times. This yields

$$f^\mathrm{I}(\mathrm{e}^{2\pi \mathrm{i}}z)=S_{11}\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+S_{12}\mathrm{\Phi }(1;0)f_\mathrm{}2(z)$$
$`(5.19)`$
with lengthy expressions for $`S_{11}`$ and $`S_{12}`$.

In a similar way, the integral representation with $`V(1;0)`$ yields

$$f^{\mathrm{II}}(z):=\underset{q=0}{}E(1;0,q)\mathrm{\Phi }(1;q)f_\mathrm{}1(z)+\mathrm{\Phi }(1;0)f_\mathrm{}2(z)$$
$`(5.20)`$
or

$$f^{\mathrm{II}}(z)=\frac{1}{\mathrm{sin}(\pi \mu (1))}\sigma _0(1)\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+\mathrm{\Phi }(1;0)f_\mathrm{}2(z)$$
$`(5.21)`$
and leads, after the loop has been traversed once, to

$$\begin{array}{cc}& \mathrm{exp}(2\pi \text{i}\lambda /3)f^{\mathrm{II}}(\mathrm{e}^{2\pi \mathrm{i}/3}z)=\{2\text{i}\mathrm{exp}(\pi \text{i}\mu (1))\sigma _0(1)+\frac{\mathrm{exp}(2\pi \text{i}\mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\hfill \\ & \frac{4\mathrm{exp}(\pi \text{i}[\mu (1)+\mu (1)])}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _1(1)\sigma _0(1)\}\mathrm{\Phi }(1;0)f_\mathrm{}1(z)\hfill \\ & +\mathrm{exp}(2\pi \text{i}\mu (1))\{1+\frac{2\text{i}\mathrm{exp}(\pi \text{i}\mu (1))}{\mathrm{sin}(\pi \mu (1))}\sigma _1(1)\sigma _1(1)\}\mathrm{\Phi }(1;0)f_\mathrm{}2(z)\hfill \end{array}$$
$`(5.22)`$
and, after three times,

$$f^{\mathrm{II}}(\mathrm{e}^{2\pi \mathrm{i}}z)=S_{21}\mathrm{\Phi }(1;0)f_\mathrm{}1(z)+S_{22}\mathrm{\Phi }(1;0)f_\mathrm{}2(z),$$
$`(5.23)`$
where again the expressions for $`S_{21}`$ and $`S_{22}`$ are too lengthy to be displayed here.

What we really want to obtain are the circuit relations for $`f_\mathrm{}1`$ and $`f_\mathrm{}2`$, that is

$$\begin{array}{cc}& f_\mathrm{}1(\mathrm{e}^{2\pi \mathrm{i}}z)=T_{11}f_\mathrm{}1(z)+T_{12}f_\mathrm{}2(z),\hfill \\ & f_\mathrm{}2(\mathrm{e}^{2\pi \mathrm{i}}z)=T_{21}f_\mathrm{}1(z)+T_{22}f_\mathrm{}2(z),\hfill \end{array}$$
$`(5.24)`$
where

$$\begin{array}{cc}& T_{11}=[S_{11}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{21}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}],\hfill \\ & T_{12}=[\mathrm{\Phi }(1;0)/\mathrm{\Phi }(1;0)][S_{12}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{22}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}],\hfill \\ & T_{21}=[\mathrm{\Phi }(1;0)/\mathrm{\Phi }(1;0)][S_{21}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{11}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}],\hfill \\ & T_{22}=[S_{22}\frac{\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))}S_{12}]/[1\frac{\sigma _0(1)\sigma _0(1)}{\mathrm{sin}(\pi \mu (1))\mathrm{sin}(\pi \mu (1))}].\hfill \end{array}$$
$`(5.25)`$
It turns out that each numerator contains a common factor which compensates the denominator. The result then is

Theorem 2. The coefficients in the circuit relations (5.24) are

$$\begin{array}{cc}& T_{11}=\mathrm{exp}(2\pi \text{i}\tau (1))\hfill \\ & +4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _0(1)\sigma _0(1)+4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _0(1)+4\sigma _2(1)\sigma _0(1)\hfill \\ & +4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)+4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _0(1)+4\sigma _1(1)\sigma _0(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _0(1)+16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)+16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +64\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1),\hfill \end{array}$$
$`(5.26a)`$
$$\begin{array}{cc}& T_{12}=(2s_0)^{(1/3)[\tau (1)\tau (1)]}\{2\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\hfill \\ & +2\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)+2\text{i}\mathrm{exp}(2\pi \text{i}\tau (1))\sigma _0(1)\hfill \\ & +8\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)+8\text{i}\sigma _2(1)\sigma _1(1)\sigma _0(1)\hfill \\ & +8\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)+8\text{i}\sigma _2(1)\sigma _1(1)\sigma _1(1)\hfill \\ & +32\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\},\hfill \end{array}$$
$`(5.26b)`$
$$\begin{array}{cc}& T_{21}=(2s_0)^{(1/3)[\tau (1)\tau (1)]}\{2\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _0(1)\hfill \\ & 2\text{i}\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)2\text{i}\sigma _1(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)\hfill \\ & 8\text{i}\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _0(1)\hfill \\ & 32\text{i}\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\},\hfill \end{array}$$
$`(5.26c)`$
$$\begin{array}{cc}& T_{22}=\mathrm{exp}(2\pi \text{i}\tau (1))\hfill \\ & +4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)+4\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)+4\sigma _2(1)\sigma _1(1)\hfill \\ & +16\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1).\hfill \end{array}$$
$`(5.26d)`$
It turns out that

$$T_{11}T_{22}T_{12}T_{21}=1.$$
$`(5.27)`$
We are looking for a multiplicative solution such that (5.2) holds. It then follows from (5.1), 5.2), (5.24) that

$$\begin{array}{cc}& (T_{11}p)\alpha +T_{21}\beta =0,\hfill \\ & T_{12}\alpha +(T_{22}p)\beta =0\hfill \end{array}$$
$`(5.28)`$
and, as a consequence, that

$$(T_{11}p)(T_{22}p)=T_{12}T_{21}$$
$`(5.29)`$
or

$$p^2(T_{11}+T_{22})p+1=0,$$
$`(5.30)`$
where (5.27) has been used. The roots $`p_1`$, $`p_2`$ of this equation satisfy

$$p_1+p_2=T_{11}+T_{22},$$
$`(5.31)`$
and they may be represented in terms of one (not necessarily real) parameter $`\omega `$ as

$$p_1=\mathrm{exp}(2\pi \text{i}\omega ),p_2=\mathrm{exp}(2\pi \text{i}\omega ).$$
$`(5.32)`$
Then we have

$$p_1+p_2=2\mathrm{cos}(2\pi \omega )=T_{11}+T_{22},$$
$`(5.33)`$
and the final result is

Theorem 3. The characteristic exponent $`\omega `$ of the multiplicative solutions is given by

$$\mathrm{cos}(2\pi \omega )=\mathrm{cos}(2\pi \tau (1))+X,$$
$`(5.34)`$
where

$$\begin{array}{cc}& X=2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _0(1)\sigma _0(1)+2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _0(1)+2\sigma _2(1)\sigma _0(1)\hfill \\ & +2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)+2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _0(1)+2\sigma _1(1)\sigma _0(1)\hfill \\ & +2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)+2\mathrm{exp}(\frac{4}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)+2\sigma _2(1)\sigma _1(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _0(1)+8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)+8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _2(1)\sigma _2(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +8\mathrm{exp}(\frac{2}{3}\pi \text{i}\tau (1))\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1)\hfill \\ & +32\sigma _2(1)\sigma _2(1)\sigma _1(1)\sigma _1(1)\sigma _0(1)\sigma _0(1).\hfill \end{array}$$
$`(5.35)`$
For each of the roots $`p_1`$ and $`p_2`$ we may determine the ratio of $`\alpha `$ and $`\beta `$ from the upper or lower equation of (5.28). Choosing in each case a convenient normalization, we get the desired multiplicative solutions

$$\begin{array}{cc}& f_{p1}(z)=(T_{22}p_1)f_\mathrm{}1(z)T_{12}f_\mathrm{}2(z),\hfill \\ & f_{p2}(z)=(T_{22}p_2)f_\mathrm{}1(z)T_{12}f_\mathrm{}2(z).\hfill \end{array}$$
$`(5.36)`$
Here the lower equation of (5.28) has been used in both cases so that $`T_{11}`$, which consists of a considerably longer expression than $`T_{22}`$, does not appear.

Whenever $`D_m=0`$ for all the $`m=1,2,\mathrm{},6`$, the origin is a regular singular point of the differential equation with exponents $`\omega =L,L`$. We are not able, however, to see analytically that then the lengthy expression (5.35) reduces to $`X=\mathrm{cos}(2\pi L)\mathrm{cos}(2\pi \tau (1))`$, but this is confirmed in examples of numerical computations, as expected.

6. Asymptotic behaviour of the late coefficients of the formal power series solutions

According to (4.12), the leading terms of the asymptotic behaviour of the $`b`$-coefficients for large $`m`$ are given by

$$\begin{array}{cc}& b(\kappa ;m,n_1,n_2)\frac{1}{\pi }(2\kappa s_0)^m\frac{\mathrm{\Gamma }(1+\mu (\kappa )+m)}{\mathrm{\Gamma }(1+m)}\hfill \\ & \times \underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+m)\hfill \\ & \times \frac{1}{(n_1p_1)!}(2\kappa t_{10})^{n_1p_1}\frac{1}{(n_2p_2)!}(2\kappa t_{20})^{n_2p_2}.\hfill \end{array}$$
$`(6.1)`$
This result may be used to discuss the asymptotic behaviour of the late $`a`$-coefficients of the formal solutions. Writing the decomposition (2.17) of the $`a`$-coefficients in terms of the $`b`$-coefficients separately for each of three consecutive indices $`n=3N,3N+1,3N+2`$, we have

$$a_{3N}(\kappa )=\underset{l=0}{\overset{2N}{}}\underset{(n_1,n_2)J_{3l}}{}b(\kappa ;N+l,n_1,n_2)=b(\kappa ;N,0,0)+\mathrm{},$$
$`(6.2)`$
$$a_{3N+1}(\kappa )=\underset{l=1}{\overset{2N+1}{}}\underset{(n_1,n_2)J_{3l1}}{}b(\kappa ;N+l,n_1,n_2)=b(\kappa ;N+1,1,0)+b(\kappa ;N+1,0,2)+\mathrm{},$$
$`(6.3)`$
$$a_{3N+2}(\kappa )=\underset{l=1}{\overset{2N+2}{}}\underset{(n_1,n_2)J_{3l2}}{}b(\kappa ;N+l,n_1,n_2)=b(\kappa ;N+1,0,1)+\mathrm{},$$
$`(6.4)`$
where

$$J_l=\{(n_1,n_2):2n_1+n_2=l\}\text{N}_0\times \text{N}_0,l=0,1,\mathrm{}.$$
$`(6.5)`$
Inserting (6.1) and omitting terms of relative order $`N^1`$, we get from (6.2)

$$\begin{array}{cc}\hfill a_{3N}(\kappa )& C_{3N}(\kappa )\underset{l=0}{\overset{2N}{}}\underset{(n_1,n_2)J_{3l}}{}\underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}p_1\frac{1}{3}p_2+l+N)}{\mathrm{\Gamma }(\mu (\kappa )+N)}\hfill \\ & \times (2\kappa s_0)^l\frac{1}{(n_1p_1)!}(2\kappa t_{10})^{n_1p_1}\frac{1}{(n_2p_2)!}(2\kappa t_{20})^{n_2p_2},\hfill \end{array}$$
$`(6.6)`$
where

$$\begin{array}{cc}\hfill C_{3N}(\kappa )& =\frac{1}{\pi }\frac{\mathrm{\Gamma }(1+\mu (\kappa )+N)\mathrm{\Gamma }(\mu (\kappa )+N)}{\mathrm{\Gamma }(1+N)}(2\kappa s_0)^N\hfill \\ & \frac{1}{\pi }\mathrm{\Gamma }(\mu (\kappa )\mu (\kappa )+N)(2\kappa s_0)^N.\hfill \end{array}$$
$`(6.7)`$
Let us now discuss (6.6) in detail: Here only such values of $`n_1`$ and $`n_2`$ occur for which

$$\frac{2}{3}n_1+\frac{1}{3}n_2=l.$$

If again terms of relative order $`N^1`$ are omitted, the ratio of the gamma functions is equal to

$$N^{(2/3)p_1(1/3)p_2+l}=N^{(2/3)(n_1p_1)+(1/3)(n_2p_2)}.$$

The asymptotic behaviour of (6.6) then becomes

$$\begin{array}{cc}& a_{3N}(\kappa )C_{3N}(\kappa )\underset{l=0}{\overset{2N}{}}\underset{(n_1,n_2)J_{3l}}{}\underset{p_1=0}{\overset{n_1}{}}\underset{p_2=0}{\overset{n_2}{}}e(\kappa ;p_1,p_2)\hfill \\ & \times \frac{1}{(n_1p_1)!}[N^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})]^{n_1p_1}\frac{1}{(n_2p_2)!}[N^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})]^{n_2p_2}.\hfill \end{array}$$
$`(6.8)`$
Here the integer power of $`2\kappa s_0`$ has been splitted in two factors with fractional powers, the meaning of which is given in terms of the phase factor $`\eta `$ of (5.10) by

$$(2\kappa s_0)^l=\{\begin{array}{c}(2s_0)^{(2/3)n_1}(2s_0)^{(1/3)n_2}\mathrm{if}\kappa =1\\ (2s_0)^{(2/3)n_1}\eta ^{2n_1}(2s_0)^{(1/3)n_2}\eta ^{n_2}\mathrm{if}\kappa =1\end{array}.$$
$`(6.9)`$
With $`n/3`$ in place of $`N`$, the terms in (6.8) look like the terms of the expansion of the exponential function

$$\begin{array}{cc}\hfill EX_0(n):& =\mathrm{exp}((\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20}))\hfill \\ & =1+(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{},\hfill \end{array}$$
$`(6.10)`$
but because of the three possible values of a third root there are two other such functions,

$$\begin{array}{cc}\hfill EX_1(n):& =\mathrm{exp}(\eta ^4(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^2(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20}))\hfill \\ & =1+\eta ^4(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^2(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{},\hfill \end{array}$$
$`(6.11)`$
$$\begin{array}{cc}\hfill EX_2(n):& =\mathrm{exp}(\eta ^2(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^4(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20}))\hfill \\ & =1+\eta ^2(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\eta ^4(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{}.\hfill \end{array}$$
$`(6.12)`$
Therefore (6.8) asymptotically shows exponential behaviour given by a certain linear combination of these three exponential functions. They can be identified by the constant term and the two linear terms of their expansion shown above. It is convenient, and easy because of the properties (5.11) of $`\eta `$, to introduce three functions which are linear combinations of the exponential functions such that only one of the identifying terms is different from zero,

$$\begin{array}{cc}& L_0(n):=\frac{1}{3}[EX_0(n)+EX_1(n)+EX_2(n)]=1+\mathrm{},\hfill \\ & L_1(n):=\frac{1}{3}[EX_0(n)+\eta ^2EX_1(n)+\eta ^4EX_2(n)]=(\frac{1}{3}n)^{2/3}(2\kappa s_0)^{2/3}(2\kappa t_{10})+\mathrm{},\hfill \\ & L_2(n):=\frac{1}{3}[EX_0(n)+\eta ^4EX_1(n)+\eta ^2EX_2(n)]=(\frac{1}{3}n)^{1/3}(2\kappa s_0)^{1/3}(2\kappa t_{20})+\mathrm{}.\hfill \end{array}$$
$`(6.13)`$
We can now determine the asymptotic behaviour of (6.8) by looking at the contributions from $`(p_1,p_2)=(n_1,n_2)`$ or $`(p_1,p_2)=(n_11,n_2)`$ or $`(p_1,p_2)=(n_1,n_21)`$, respectively. This yields

$$a_{3N}(\kappa )C_{3N}(\kappa )\left\{S_0(\kappa )L_0(3N)+(2\kappa s_0)^{1/3}S_1(\kappa )L_1(3N)+(2\kappa s_0)^{2/3}S_2(\kappa )L_2(3N)\right\}.$$
$`(6.14)`$
In a similar way we may obtain from (6.3) and (6.4), respectively,

$$\begin{array}{cc}& a_{3N+1}(\kappa )C_{3N+1}(\kappa )\hfill \\ & \times \left\{(2\kappa s_0)^{2/3}S_2(\kappa )L_0(3N+1)+S_0(\kappa )L_1(3N+1)+(2\kappa s_0)^{1/3}S_1(\kappa )L_2(3N+1)\right\},\hfill \end{array}$$
$`(6.15)`$
$$\begin{array}{cc}& a_{3N+2}(\kappa )C_{3N+2}(\kappa )\hfill \\ & \times \left\{(2\kappa s_0)^{1/3}S_1(\kappa )L_0(3N+2)+(2\kappa s_0)^{2/3}S_2(\kappa )L_1(3N+2)+S_0(\kappa )L_2(3N+2)\right\}.\hfill \end{array}$$
$`(6.16)`$
If we now switch back to a representation in terms of the exponential functions (6.10)-(6.12), Stokes multipliers, according to (5.13) above, appear as their factors,

$$\begin{array}{cc}& a_{3N}(\kappa )\frac{1}{3}C_{3N}(\kappa )\hfill \\ & \times \left\{\sigma _0(\kappa )EX_0(3N)+\sigma _1(\kappa )EX_1(3N)+\sigma _2(\kappa )EX_2(3N)\right\},\hfill \\ & a_{3N+1}(\kappa )\frac{1}{3}C_{3N+1}(\kappa )\hfill \\ & \times \left\{\sigma _0(\kappa )EX_0(3N+1)+\eta ^2\sigma _1(\kappa )EX_1(3N+1)+\eta ^4\sigma _2(\kappa )EX_2(3N+1)\right\},\hfill \\ & a_{3N+2}(\kappa )\frac{1}{3}C_{3N+2}(\kappa )\hfill \\ & \times \left\{\sigma _0(\kappa )EX_0(3N+2)+\eta ^4\sigma _1(\kappa )EX_1(3N+2)+\eta ^2\sigma _2(\kappa )EX_2(3N+2)\right\}.\hfill \end{array}$$
$`(6.17)`$
These three asymptotic equations can be combined to give

Theorem 4. The leading terms of the asymptotic exponential behaviour, as $`n\mathrm{}`$, of the coefficients of the formal solutions are given by

$$\begin{array}{cc}& a_n(\kappa )\frac{1}{3}\pi ^1(2\kappa s_0)^{n/3}\mathrm{\Gamma }(\frac{1}{3}[\tau (\kappa )\tau (\kappa )+n])\hfill \\ & \times [\sigma _0(\kappa )EX_0(n)+\eta ^{2n}\sigma _1(\kappa )EX_1(n)+\eta ^{4n}\sigma _2(\kappa )EX_2(n)].\hfill \end{array}$$
$`(6.18)`$
This is a concrete example, with all the quantities determined explicitly, of the structural results obtained by Immink , and it is interesting also in the context of related work by other authors .

7. Postponed proofs

7.1 Proof of Lemma 1

Let

$$f(z)=z^\lambda u(z)$$
$`(7.1)`$
and multiply the differential equation for $`u(z)`$ by $`z^6`$ in order to remove all the negative powers of $`z`$. We then are concerned with the differential equation

$$\begin{array}{cc}\hfill L_zu(z):=& z^8u^{\prime \prime }+(2\lambda +1)z^7u^{}\hfill \\ & [\underset{m=1}{\overset{6}{}}D_mz^{6m}+(L^2\lambda ^2)z^6+\underset{m=1}{\overset{6}{}}B_mz^{6+m}]u(z)=0.\hfill \end{array}$$
$`(7.2)`$
We are looking for a solution of this differential equation in the form of an integral representation

$$u(z)=_{C_{t_2}}_{C_{t_1}}_{C_s}K(z;s,t_1,t_2)v(s,t_1,t_2)\text{d}s\text{d}t_1\text{d}t_2$$
$`(7.3)`$
with the kernel

$$K=K(z;s,t_1,t_2)=\mathrm{exp}(z^3s+z^2t_1+zt_2).$$
$`(7.4)`$
If we perform the differentiations with respect to $`z`$ under the integrals, the differential equation becomes

$$L_zu=vL_zK\text{d}s\text{d}t_1\text{d}t_2=0,$$
$`(7.5)`$
where

$$\begin{array}{cc}\hfill \frac{1}{K}L_zK& =(9s^2B_6)z^{12}+(12st_1B_5)z^{11}+(6st_2+4t_1^2B_4)z^{10}\hfill \\ & +([6\lambda +9]s+4t_1t_2B_3)z^9+([4\lambda +4]t_1+t_2^2B_2)z^8\hfill \\ & +([2\lambda +1]t_2B_1)z^7+(\lambda ^2L^2)z^6\underset{m=1}{\overset{6}{}}D_mz^{6m}.\hfill \end{array}$$
$`(7.6)`$
It is advisable to rewrite this in terms of the quantities $`p_1`$, $`p_2`$, $`p_3`$, according to (1.5), which determine the exponential factors of the formal solutions (1.2)-(1.3), or in terms of the related quantities $`s_0`$, $`t_{10}`$, $`t_{20}`$ according to (2.3). Then (7.6) becomes

$$\begin{array}{cc}\hfill \frac{1}{K}L_zK& =9(s^2s_0^2)z^{12}+12(st_1s_0t_{10})z^{11}+(6[st_2s_0t_{20}]+4[t_1^2t_{10}^2])z^{10}\hfill \\ & +([6\lambda +9]s+4[t_1t_2t_{10}t_{20}]+[4t_{10}t_{20}B_3])z^9\hfill \\ & +([4\lambda +4]t_1+[t_2^2t_{20}^2]+[t_{20}^2B_2])z^8\hfill \\ & +([2\lambda +1]t_2B_1)z^7+(\lambda ^2L^2)z^6\underset{m=1}{\overset{6}{}}D_mz^{6m}.\hfill \end{array}$$
$`(7.7)`$
By repeated partial integrations of the exponential function with respect to $`s`$ or $`t_1`$ or $`t_2`$, respectively, it is possible to get rid of all the powers of $`z`$. This lengthy task can more conveniently be performed in a formal way as follows: Let us find a partial differential expression $`M=M_{s,t_1,t_2}`$ with respect to $`s`$, $`t_1`$, $`t_2`$, independent of $`z`$, such that

$$L_zKM_{s,t_1,t_2}K.$$
$`(7.8)`$
The powers of $`z`$ correspond to partial derivatives with respect to $`s`$ or $`t_1`$ or $`t_2`$, and there seems to be some ambiguity as to the choice of the partial derivatives which yield the same total power of $`z`$. This ambiguity is resolved by the following reasonning: In order to get the order of each derivative as small as possible, we would prefer derivatives with respect to $`s`$, which account for $`z^3`$, with highest priority, next $`t_1`$, which accounts for $`z^2`$, last $`t_2`$, which accounts only for $`z`$. Conflicting with this policy, however, are some other requirements which ensure that we get a differential expression appropriate for our purpose. So the singularities should be at the right places, which are related to the coefficients of the exponential factor of the formal solutions we want to represent. In particular, bilinear or quadratic factors, such as $`t_1^2`$ for instance, should occur in the form $`t_1^2t_{10}^2`$. For this reason we have in (7.7) already added and subtracted the term $`4t_{10}^2`$. It is then necessary to treat these two terms differently, so that any power is always multiplied by the corresponding derivative. In the example just mentioned this means that $`z^{10}`$ has to be translated into $`^4/s^3t_2`$ for one term, but into $`^4/s^2t_1^2`$ for the other. The unique result now is

$$\begin{array}{cc}\hfill M_{s,t_1,t_2}& =9(s^2s_0^2)^4/s^4+12(st_1s_0t_{10})^4/s^3t_1\hfill \\ & +6(st_2s_0t_{20})^4/s^3t_2+4(t_1^2t_{10}^2)^4/s^2t_1^2\hfill \\ & +([6\lambda +9]s+[4t_{10}t_{20}B_3])^3/s^3+4[t_1t_2t_{10}t_{20}]^4/s^2t_1t_2\hfill \\ & +([4\lambda +4]t_1+[t_{20}^2B_2])^3/s^2t_1+[t_2^2t_{20}^2]^4/s^2t_2^2\hfill \\ & +([2\lambda +1]t_2B_1)^3/s^2t_2+(\lambda ^2L^2)^2/s^2\hfill \\ & D_1^2/st_1D_2^2/st_2D_3/sD_4/t_1D_5/t_2D_6.\hfill \end{array}$$
$`(7.9)`$
Next let us introduce the adjoint differential expression $`\overline{M}=\overline{M}_{s,t_1,t_2}`$ defined, with any sufficiently differentiable function $`v=v(s,t_1,t_2)`$, by

$$\begin{array}{cc}& \overline{M}_{s,t_1,t_2}v=(^4/s^4)\{9(s^2s_0^2)v\}+(^4/s^3t_1)\{12(st_1s_0t_{10})v\}\hfill \\ & +(^4/s^3t_2)\{6(st_2s_0t_{20})v\}+(^4/s^2t_1^2)\{4(t_1^2t_{10}^2)v\}\hfill \\ & (^3/s^3)\{([6\lambda +9]s+[4t_{10}t_{20}B_3])v\}+(^4/s^2t_1t_2)\{4(t_1t_2t_{10}t_{20})v\}\hfill \\ & (^3/s^2t_1)\{([4\lambda +4]t_1+[t_{20}^2B_2])v\}+(^4/s^2t_2^2)\{(t_2^2t_{20}^2)v\}\hfill \\ & (^3/s^2t_2)\{([2\lambda +1]t_2B_1)v\}+(\lambda ^2L^2)(^2v/s^2)\hfill \\ & D_1(^2v/st_1)D_2(^2v/st_2)+D_3(v/s)+D_4(v/t_1)+D_5(v/t_2)D_6v.\hfill \end{array}$$
$`(7.10)`$
The difference as compared with $`M`$ is that the factors in front of each derivative have here to be differentiated too and that all the terms of odd order change their sign. The usefulness of the adjoint expression lies in the formula, known in the one-variable case as the identity of Lagrange,

$$vMKK\overline{M}v=RHS,$$
$`(7.11)`$
where the right-hand side $`RHS`$ is a lengthy expression, bilinear in $`v`$ and $`K`$ or their partial derivatives, which may be so arranged that it consists of a sum of terms of which each is a total derivative with respect to one of the variables. For the term with the factor $`D_1`$, for instance, this reads

$$v\frac{^2K}{st_1}K\frac{^2v}{st_1}=\frac{1}{2}\frac{}{s}[v\frac{K}{t_1}K\frac{v}{t_1}]+\frac{1}{2}\frac{}{t_1}[v\frac{K}{s}K\frac{v}{s}].$$
$`(7.12)`$
By means of (7.5) and (7.8), the equation to be satisfied now is

$$\begin{array}{cc}\hfill L_zu=& vM_{s,t_1,t_2}K\text{d}s\text{d}t_1\text{d}t_2\hfill \\ \hfill =& K\overline{M}_{s,t_1,t_2}v\text{d}s\text{d}t_1\text{d}t_2+RHS\text{d}s\text{d}t_1\text{d}t_2=0.\hfill \end{array}$$
$`(7.13)`$
The first term in the second line can be made to vanish if $`v`$ is required to be a solution of the partial differential equation

$$\overline{M}_{s,t_1,t_2}v(s,t_1,t_2)=0,$$
$`(7.14)`$
and the second term by a suitable choice of the contours of integration, for it is a sum of semi-integrated terms, each involving the difference of the values of the integrand at the termini of the contour of one variable and only two remaining integrals with respect to the other two variables. The partial differential equation (2.2) is the same as (7.14), after the derivatives of the products in (7.10) have been resolved. This completes the proof of Lemma 1.

7.2 Proof of Lemma 2

In terms of the shifted variables

$$S=s\kappa s_0,T_1=t_1\kappa t_{10},T_2=t_2\kappa t_{20},$$
$`(7.15)`$
the differential equation (2.2) reads

$$\begin{array}{cc}& 9S(S+2\kappa s_0)(^4v/S^4)+12(ST_1+\kappa t_{10}S+\kappa s_0T_1)(^4v/S^3T_1)\hfill \\ & +6(ST_2+\kappa t_{20}S+\kappa s_0T_2)(^4v/S^3T_2)+4T_1(T_1+2\kappa t_{10})(^4v/S^2T_1^2)\hfill \\ & +T_2(T_2+2\kappa t_{20})(^4v/S^2T_2^2)+4(T_1T_2+\kappa t_{20}T_1+\kappa t_{10}T_2)(^4v/S^2T_1T_2)\hfill \\ & +([816\lambda ]\kappa s_0+B_34t_{10}t_{20})(^3v/S^3)+(816\lambda )S(^3v/S^3)\hfill \\ & +([524\lambda ]\kappa t_{10}+B_2t_{20}^2)(^3v/S^2T_1)+(524\lambda )T_1(^3v/S^2T_1)\hfill \\ & +([252\lambda ]\kappa t_{20}+B_1)(^3v/S^2T_2)+(252\lambda )T_2(^3v/S^2T_2)\hfill \\ & +([\lambda 12]^2L^2)(^2v/S^2)D_1(^2v/ST_1)D_2(^2v/ST_2)\hfill \\ & +D_3(v/S)+D_4(v/T_1)+D_5(v/T_2)D_6v=0.\hfill \end{array}$$
$`(7.16)`$
Inserting a power series solution

$$v=\underset{m=0}{}\underset{n_1=0}{}\underset{n_2=0}{}a(m,n_1,n_2)S^{\mu +m}T_1^{\nu _1n_1}T_2^{\nu _2n_2},$$
$`(7.17)`$
we obtain for the coefficients the recurrence relation

$$\begin{array}{cc}& 6\kappa s_0(\mu +m)(\mu +m1)(\mu +m2)[3(\mu +m)2(\nu _1+n_1)(\nu _2+n_2)\lambda \tau (\kappa )+6]\hfill \\ & \times a(m,n_1,n_2)\hfill \\ & +(\mu +m1)(\mu +m2)\{[3(\mu +m)2(\nu _1+n_1)(\nu _2+n_2)\lambda +3]^2L^2\}\hfill \\ & \times a(m1,n_1,n_2)\hfill \\ & +(\mu +m1)(\mu +m2)(\nu _1+n_11)\{4\kappa t_{10}[3(\mu +m)+2(\nu _1+n_1)+(\nu _2+n_2)\hfill \\ & +\lambda 4]+t_{20}^2B_2\}a(m1,n_11,n_2)\hfill \\ & +(\mu +m1)(\mu +m2)(\nu _2+n_21)\{2\kappa t_{20}[3(\mu +m)+2(\nu _1+n_1)+(\nu _2+n_2)\hfill \\ & +\lambda \frac{7}{2}]+B_1\}a(m1,n_1,n_21)\hfill \\ & +D_1(\mu +m2)(\nu _1+n_11)a(m2,n_11,n_2)\hfill \\ & +D_2(\mu +m2)(\nu _2+n_21)a(m2,n_1,n_21)\hfill \\ & +D_3(\mu +m2)a(m2,n_1,n_2)D_4(\nu _1+n_11)a(m3,n_11,n_2)\hfill \\ & D_5(\nu _2+n_21)a(m3,n_1,n_21)D_6a(m3,n_1,n_2)=0,\hfill \end{array}$$
$`(7.18)`$
valid for $`m0`$, $`n_10`$, $`n_20`$ provided that we agree that all the $`a`$-coefficients are equal to zero if any of the indices $`m,n_1,n_2`$ is less than zero. Assuming that $`a(0,0,0)0`$, we get from the equation for $`m=n_1=n_2=0`$ the indicial equation

$$\mu (\mu 1)(\mu 2)(3\mu 2\nu _1\nu _2\lambda \tau (\kappa )+6)=0,$$
$`(7.19)`$
with $`\tau (\kappa )`$ according to (1.6). Possible values of the exponent $`\mu `$ are therefore $`\mu =0,1,2,`$ or

$$\mu =\frac{1}{3}[2\nu _1+\nu _2+\lambda +\tau (\kappa )]2.$$
$`(7.20)`$
In order to avoid complications, we may assume that the last possibility does not yield an integer value. This can always be guarantied be a suitable choice of the still disposable parameter $`\lambda `$. Inspection of the recurrence relation then shows that each of the possible values of $`\mu `$ leads to a solution of the partial differential equation. We here need not further consider the three solutions which are regular with respect to $`S`$ at $`S=0`$, but we continue to discuss the singular one with the exponent (7.20), writing for the associated coefficients $`A(\kappa ;m,n_1,n_2)`$ rather than $`a(m,n_1,n_2)`$. Simplifying by means of (7.20) and introducing the $`b`$-coefficients according to (2.6) , we get (2.7). The coefficients for which $`3m2n_1n_2=0`$ are constants of integration and may be chosen to be zero. This completes the proof of Lemma 2.

7.3 Proof of the limit formula

We have to verify that the limit formula (4.14) satisfies (4.13 ). Substituting it for the $`e`$-coefficients in (4.13) and interchanging the summations, we have to evaluate sums such as

$$\underset{p=j}{\overset{n}{}}\frac{(1)^{np}}{(pj)!(np)!}=\frac{(1)^{nj}}{(nj)!}\underset{q=0}{\overset{nj}{}}(1)^q\left(\genfrac{}{}{0pt}{}{nj}{q}\right)=\{\begin{array}{c}1\mathrm{if}j=n\\ 0\mathrm{if}0jn1,n>0\end{array}.$$
$`(7.21)`$
Therefore only one term survives on the right-hand side, which then becomes equal to the left.

7.4 Proof of Theorem 1

In order to prove Theorem 1, we first multiply the continuation formula (3.5), with $`r_1=r_2=0`$, by

$$(\frac{1}{2}\frac{s}{2\kappa s_0})^{\mu (\kappa )}=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(\mu (\kappa ))_j}{j!}(\frac{1}{2}+\frac{s}{2\kappa s_0})^j,$$
$`(7.22)`$
where the left-hand side is used on the left and the right-hand side on the right of the continuation formula. Then the left-hand side is the same as above (4.1) with the power factor in front of the series removed, but the right becomes, after multiplication of the two power series,

$$\begin{array}{cc}& \underset{q_1=0}{\overset{\mathrm{}}{}}\underset{q_2=0}{\overset{\mathrm{}}{}}E(\kappa ;0,0;q_1,q_2)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2)\hfill \\ & \times \underset{k=0}{\overset{\mathrm{}}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\mathrm{\Gamma }(\nu _1+q_1+l_1)\mathrm{\Gamma }(\nu _2+q_2+l_2)H(\kappa ;k,l_1,l_2;q_1,q_2)\hfill \\ & \times (\frac{1}{2}+\frac{s}{2\kappa s_0})^{\mu (\kappa )+{\scriptscriptstyle \frac{2}{3}}q_1+{\scriptscriptstyle \frac{1}{3}}q_2+k}(t_1+\kappa t_{10})^{\nu _1q_1l_1}(t_2+\kappa t_{20})^{\nu _2q_2l_2},\hfill \end{array}$$
$`(7.23)`$
with $`H`$ defined in (4.16). Proceeding as above, and making use of the formula

$$(\alpha k)_m=(\alpha )_m\frac{(1\alpha )_k}{(1\alpha m)_k},$$
$`(7.24)`$
we obtain

$$\begin{array}{cc}& \frac{\pi }{\mathrm{sin}(\pi \mu (\kappa ))}\frac{1}{\mathrm{\Gamma }(1+\mu (\kappa )+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{q_1=0}{\overset{n_1}{}}\underset{q_2=0}{\overset{n_2}{}}E(\kappa ;0,0;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{m!}\hfill \\ & \times \underset{k=0}{\overset{\mathrm{}}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_k}{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2m)_k}H(\kappa ;k,l_1,l_2;q_1,q_2)\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1+q_1+l_1=n_1}\left(\frac{(1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2+q_2+l_2=n_2}.\hfill \end{array}$$
$`(7.25)`$
Keeping the first $`K+1`$ singular terms on the right and solving for the $`E`$-coefficient with $`q_1=n_1`$, $`q_2=n_2`$, we have

$$\begin{array}{cc}& E(\kappa ;0,0;n_1,n_2)\hfill \\ & \times [1+\underset{k=1}{\overset{K}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2)_k}{(1+\mu (\kappa )+\frac{2}{3}n_1+\frac{1}{3}n_2m)_k}H(\kappa ;k,0,0;n_1,n_2)+O(m^{K1})]\hfill \\ & =\frac{\pi }{\mathrm{sin}(\pi \mu (\kappa ))}\frac{m!}{\mathrm{\Gamma }(1+\mu (\kappa )+m)\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}(2\kappa s_0)^mb(\kappa ;m,n_1,n_2)\hfill \\ & \underset{(q_1,q_2)(n_1,n_2)}{\underset{q_1=0}{\overset{n_1}{}}\underset{q_2=0}{\overset{n_2}{}}}E(\kappa ;0,0;q_1,q_2)\frac{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}q_1\frac{1}{3}q_2+m)}{\mathrm{\Gamma }(\mu (\kappa )\frac{2}{3}n_1\frac{1}{3}n_2+m)}\hfill \\ & \times [1+\underset{k=1}{\overset{K}{}}\underset{l_1=0}{\overset{k}{}}\underset{l_2=0}{\overset{k}{}}\frac{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2)_k}{(1+\mu (\kappa )+\frac{2}{3}q_1+\frac{1}{3}q_2m)_k}H(\kappa ;k,l_1,l_2;q_1,q_2)+O(m^{K1})]\hfill \\ & \times \left(\frac{1}{j_1!}(2\kappa t_{10})^{j_1}\right)_{j_1+q_1+l_1=n_1}\left(\frac{1}{j_2!}(2\kappa t_{20})^{j_2}\right)_{j_2+q_2+l_2=n_2}.\hfill \end{array}$$
$`(7.26)`$
This is essentially (4.15) because of (4.11) and completes the proof of Theorem 1.

7.5 Choice of the computational parameter $`\lambda `$

If $`D_3=D_6=0`$, the recurrence relation (2.7) for $`b(\kappa ;m,0,0)`$ reduces to a two-term relation, and we obtain

$$b(\kappa ;m,0,0)=(2\kappa s_0)^m\frac{(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_m(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_m}{m!}.$$
$`(7.27)`$
Rewriting (4.16) by means of the identity

$$\frac{(x)_{kj}}{(kj)!}=\frac{(x)_k}{k!}\frac{(k)_j}{(1xk)_j},$$
$`(7.28)`$
we then have

$$H(\kappa ;k,0,0;0,0)=\frac{(\mu (\kappa ))_k}{k!}\underset{j=0}{\overset{k}{}}\frac{(k)_j(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_j(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))_j}{(1\mu (\kappa )k)_j(1+\mu (\kappa ))_jj!}.$$
$`(7.29)`$
Because of $`3\mu (\kappa )=\lambda +\tau (\kappa )3`$ according to (2.5) and (2.11) and $`\tau (\kappa )=3\tau (\kappa )`$ according to (1.6), the series is a terminating one-balanced hypergeometric series at unit argument, which can be summed by the theorem of Saalschütz , so that

$$H(\kappa ;k,0,0;0,0)=\frac{(\frac{1}{3}L+\frac{1}{3}\lambda )_k(\frac{1}{3}L+\frac{1}{3}\lambda )_k}{(1+\mu (\kappa ))_kk!}.$$
$`(7.30)`$
In total, we get

$$\begin{array}{cc}& e(\kappa ;0,0)=\frac{\pi }{\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa ))}\frac{\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa )+m)\mathrm{\Gamma }(\frac{1}{3}L+\frac{1}{3}\tau (\kappa )+m)}{\mathrm{\Gamma }(\frac{1}{3}\lambda +\frac{1}{3}\tau (\kappa )+m)\mathrm{\Gamma }(\frac{1}{3}\lambda +\frac{1}{3}\tau (\kappa )+m)}\hfill \\ & \times \left[1+\underset{k=1}{\overset{K}{}}\frac{(\frac{1}{3}L+\frac{1}{3}\lambda )_k(\frac{1}{3}L+\frac{1}{3}\lambda )_k}{(\frac{1}{3}\lambda +\frac{1}{3}\tau (\kappa )m)_kk!}+O(m^{K1})\right]^1.\hfill \end{array}$$
$`(7.31)`$
For $`\lambda =L`$ or $`\lambda =L`$, the terms with $`k=1,2,\mathrm{}`$ of the asymptotic series all vanish and the factor in front of the series becomes independent of $`m`$. More generally this means that some terms which might become quite large when L is not small can be removed by such a choice of $`\lambda `$ from the asymptotic series and incorporated in the $`m`$-dependence of the function in front of the series.

References

 W. Balser, W. B. Jurkat, and D. A. Lutz, Transfer of connection problems for meromorphic differential equations of rank $`2`$ and representations of solutions, J. Math. Anal. Appl. 85 (1982) 488-542.

 B. L. J. Braaksma, Multisummability and Stokes multipliers of linear meromorphic differential equations, J. Differ. Equations 92 (1991) 45-75.

 W. Bühring, The characteristic exponent of second-order linear differential equations with two irregular singular points, Proc. Amer. Math. Soc. 118 (1993) 801-812.

 A. Duval, Triconfluent Heun’s equation, in: A. Ronveaux, ed., Heun’s differential equations (Oxford University Press, Oxford, New York, Tokyo, 1995).

 R. Hoeppner and R. Schäfke, On the remainders of asymptotic expansions of solutions of linear differential equations near irregular singular points of higher rank, Math. Nachr. 205 (1999) 89-113.

 G. K. Immink, A note on the relationship between Stokes multipliers and formal solutions of analytic differential equations, SIAM J. Math. Anal. 21 (1990) 782-792.

 E. L. Ince, Ordinary differential equations ( Dover, New York, 1956).

 M. Kohno, Stokes phenomenon for general linear ordinary differential equations with two singular points, in: Japan - United States seminar on ordinary differential and functional equations, Lecture Notes in Mathematics 243 (Springer, 1971) 310-314.

 M. Kohno, A two point connection problem for general linear ordinary differential equations, Hiroshima Math. J. 4 (1974) 293-338.

 M. Loday-Richaud, Calcul des invariants de Birkhoff des systèmes d’ordre deux, Funk. Ekvac. 33 (1990) 161-225.

 Y. L. Luke, The special functions and their approximations, Vol. 1 (Academic Press, New York, 1969).

 J. Martinet and J.-P. Ramis, Elementary acceleration and multisummability I, Ann. Inst. Henri Poincaré, Phys. Théor. 54 (1991) 331-401.

 F. Naundorf, A connection problem for second order linear differential equations with two irregular singular points, SIAM J. Math. Anal. 7 (1976) 157-175.

 F. Naundorf, Ein Verfahren zur Berechnung der charakteristischen Exponenten von linearen Differentialgleichungen zweiter Ordnung mit zwei stark singulären Stellen, Z. Angew. Math. Mech. 57 (1977) 47-49.

 A. B. Olde Daalhuis and F. W. J. Olver, On the calculation of Stokes multipliers for linear differential equations of the second order, Methods Appl. Anal. 2 (1995) 348-367.

 F. W. J. Olver, Asymptotics and special functions (Academic Press, New York 1974).

 R. Schäfke and D. Schmidt, The connection problem for general linear ordinary differential equations at two regular singular points with applications in the theory of special functions, SIAM J. Math. Anal. 11 (1980) 848-862.

 Y. Sibuya, Global theory of a second-order linear ordinary differential equation with a polynomial coefficient (North-Holland, Amsterdam, New York, 1975).

 V. S. Varadarajan, Linear meromorphic differential equations: a modern point of view, Bull. Amer. Math. Soc. 33 (1996) 1-42.

 E. Wagenführer, Die Determinantenmethode zur Berechnung des charakteristischen Exponenten der endlichen Hillschen Differentialgleichung, Numer. Math. 35 (1980) 405-420.

 E. Wagenführer and H. Lang, Berechnung des charakteristischen Exponenten der endlichen Hillschen Differentialgleichung durch numerische Integration, Numer. Math. 32 (1979) 31-50.

 E. T. Whittaker and G. N. Watson, A course of modern analysis (Cambridge University Press, London , 1927).

Wolfgang Bühring

Physikalisches Institut

Universität Heidelberg

Philosophenweg 12

69120 Heidelberg

GERMANY

buehring@physi.uni-heidelberg.de


format: text

text:
# DATALOG with constraints — an answer-set programming system

## Introduction

Many important computational problems in combinatorial optimization, constraint satisfaction and artificial intelligence can be cast as search problems. Answer-set programming (ASP) (??) was recently identified as a declarative programming paradigm appropriate for such applications. Logic programming with the stable-model semantics (stable logic programming, for short) was proposed as an embodiment of this paradigm. Disjunctive logic programming with the answer-set semantics is another implementation of ASP currently under development (?). Early experimental results demonstrate the potential of answer-set programming approaches in such areas as planning and constraint satisfaction (???).

In this paper we describe another formalism that implements the ASP approach. We call it DATALOG with constraints and denote by $`\mathrm{DC}`$. Our goal is to design an ASP system with a semantics more readily understandable than the semantics of stable models. We seek a semantics that would be as close as possible to propositional satisfiability yet as expressive and as effective, especially from the point of view of conciseness of representations and time performance, as the stable logic programming. We argue that $`\mathrm{DC}`$ has a potential to become a practical declarative programming tool. We show that it yields intuitive and small-size encodings, we characterize its complexity and expressive power and present computational experiments demonstrating its effectiveness.

Answer-set programming is a paradigm in which programs are built as theories in some formal system $``$ with a well-defined syntax, and with a semantics that assigns to a theory $`P`$ in the system a collection of subsets of some domain. These subsets are referred to as answer sets of $`P`$ and specify the results of computation based on $`P`$. To solve a problem $`\mathrm{\Pi }`$ in an ASP formalism, we find a program $`P`$ so that the solutions to $`\mathrm{\Pi }`$ can be reconstructed, in polynomial (ideally, linear) time, from the answer sets to $`P`$.

The definition of the answer-set programming given above is very general. Essentially any logic formalism can be a basis for an answer-set programming system. For instance, the propositional logic gives rise to an ASP system: programs are collections of propositional clauses, their models are answer sets. To solve, say, a planning problem, we encode the constraints of the problem as propositional clauses in such a way that legal plans are determined by models of the resulting propositional theory. This approach, called satisfiability planning, received significant attention lately and was shown to be quite effective (???).

Recently, several implementations of the ASP approach were developed that are based on nonmonotonic logics such as smodels (?), for stable logic programming, dlv (?), for disjunctive logic programming with answer-set semantics, and deres (?), for default logic with Reiter’s extensions. All these systems have been extensively studied. Promising experimental results concerning their performance were reported (???).

The question arises which formal logics are appropriate as bases of answer-set programming implementations. To discuss such a general question one needs to formulate quality criteria with respect to which ASP systems can be compared. At the very least, these criteria should include:
1. expressive power
2. time performance
3. simplicity of the semantics
4. ease of coding, conciseness of programs.

We will discuss these criteria in detail elsewhere. We will make here only a few brief comments on the matter. From the point of view of the expressive power all the systems that we discussed are quite similar. Propositional logic and stable logic programming are well-attuned to the class NP (?). Disjunctive logic programming and default logic capture the class $`\mathrm{\Sigma }_P^2`$ (??). However, this distinction is not essential as recently pointed out in (?). The issue of time performance can be resolved only through comprehensive experimentation and this work is currently under way.

As concerns inherent complexity of the system and intuitiveness of the semantics, ASP systems based on the propositional logic seem to be clear winners. However, propositional logic is monotone and modeling indefinite information and phenomena such as the frame problem is not quite straightforward. In applications involving the computation of transitive closures, as in the problem of existence of hamilton cycles, it leads to programs that are large and, thus, difficult to process. In this respect, ASP systems based on nonmonotonic logics have an edge. They were designed to handle incomplete and indefinite information. Thus, they often yield more concise programs. However, they require more elaborate formal machinery and their semantics are more complex.

Searching for the middle ground between systems such as logic programming with stable model semantics and propositional logic, we propose here a new ASP formalism, $`\mathrm{DC}`$. Our guiding principle was to design a system which would lead to small-size encodings, believing that small theories will lead to more efficient solutions. We show that $`\mathrm{DC}`$ is nonmonotonic, has the same expressive power as stable logic programming but that its semantics stays closer to that of propositional logic. Thus, it is arguably simpler than the stable-model semantics. We present experimental results that demonstrate that $`\mathrm{DC}`$ is competitive with ASP implementations based on nonmonotonic logics (we use smodels for comparison) and those based on propositional logics (we use csat (?) in our experiments). Our results strongly indicate that formalisms which provide smaller-size encodings are more effective as practical search-problem solvers.

## DATALOG with constraints

A $`\mathrm{DC}`$ theory (or program) consists of constraints and Horn rules (DATALOG program). This fact motivates out choice of terminology — DATALOG with constraints. We start a discussion of $`\mathrm{DC}`$ with the propositional case. Our language is determined by a set of atoms $`At`$. We will assume that $`At`$ is of the form $`At=At_CAt_H`$, where $`At_C`$ and $`At_H`$ are disjoint.

A DC theory (or DC program) is a triple $`T=(T_C,T_H,T_{PC})`$, where
1. $`T_C`$ is a set of propositional clauses $`\neg a_1\mathrm{}\neg a_mb_1\mathrm{}b_n`$ such that all $`a_i`$ and $`b_j`$ are from $`At_C`$,
2. $`T_H`$ is a set of Horn rules $`a_1\mathrm{}a_mb`$ such that $`bAt_H`$ and all $`a_i`$ are from $`At`$,
3. $`T_{PC}`$ is a set of clauses over $`At`$.
By $`At(T)`$, $`At_C(T)`$ and $`At_{PC}(T)`$ we denote the set of atoms from $`At`$, $`At_C`$ and $`At_{PC}`$, respectively, that actually appear in $`T`$.

With a $`\mathrm{DC}`$ theory $`T=(T_C,T_H,T_{PC})`$ we associate a family of subsets of $`At_C(T)`$. We say that a set $`MAt_C(T)`$ satisfies $`T`$ (is an answer set of $`T`$) if
1. $`M`$ satisfies all the clauses in $`T_C`$, and
2. the closure of $`M`$ under the Horn rules in $`T_H`$, $`M^c=LM(T_HM)`$ satisfies all clauses in $`T_{PC}`$ ($`LM(P)`$ denotes the least model of a Horn program $`P`$).
Intuitively, the collection of clauses in $`T_C`$ can be thought of as a representation of the constraints of the problem, Horn rules in $`T_H`$ can be viewed as a mechanism to compute closures of sets of atoms satisfying the constraints in $`T_C`$, and the clauses in $`T_{PC}`$ can be regarded as constraints on closed sets (we refer to them as post-constraints). A set of atoms $`MAt_C(T)`$ is a model if it (propositionally) satisfies the constraints in $`T_C`$ and if its closure (propositionally) satisfies the constraints in $`T_{PC}`$. Thus, the semantics of $`\mathrm{DC}`$ retains much of the simplicity of the semantics of propositional logic.

$`\mathrm{DC}`$ can be used as a computational tool to solve search problems. We define a search problem $`\mathrm{\Pi }`$ to be determined by a set of finite instances, $`D_\mathrm{\Pi }`$, such that for each instance $`ID_\mathrm{\Pi }`$, there is a finite set $`S_\mathrm{\Pi }(I)`$ of all solutions to $`\mathrm{\Pi }`$ for the instance $`I`$. For example, the problem of finding a hamilton cycle in a graph is a search problem: graphs are instances and for each graph, its hamilton cycles (sets of their edges) are solutions. A $`\mathrm{DC}`$ theory $`T=(T_C,T_H,T_{PC})`$ solves a search problem $`\mathrm{\Pi }`$ if solutions to $`\mathrm{\Pi }`$ can be computed (in polynomial time) from answer sets to $`T`$. Propositional logic and stable logic programming are used as problem solving formalisms following the same general paradigm. To illustrate all the concepts introduced here and show how $`\mathrm{DC}`$ programs can be built by modeling problem constraints, we will now present a $`\mathrm{DC}`$ program that solves the hamilton-cycle problem.

Consider a directed graph $`G`$ with the vertex set $`V`$ and the edge set $`E`$. Consider a set of atoms $`\{hc(a,b):(a,b)E\}`$. An intuitive interpretation of an atom $`hc(a,b)`$ is that the edge $`(a,b)`$ is in a hamilton cycle. Include in $`T_C`$ all clauses of the form $`\neg hc(b,a)\neg hc(c,a)`$, where $`a,b,cV`$, $`bc`$ and $`(b,a),(c,a)E`$. In addition, include in $`T_C`$ all clauses of the form $`\neg hc(a,b)\neg hc(a,c)`$, where $`a,b,cV`$, $`bc`$ and $`(a,b),(a,c)E`$. Clearly, the set of propositional variables of the form $`\{hc(a,b):(a,b)F\}`$, where $`FE`$, satisfies all clauses in $`T_C`$ if and only if no two distinct edges in $`F`$ end in the same vertex and no two distinct edges in $`F`$ start in the same vertex. In other words, $`F`$ spans a collection of paths and cycles in $`G`$.

To guarantee that the edges in $`F`$ define a hamilton cycle, we must enforce that all vertices of $`G`$ are reached by means of the edges in $`F`$ if we start in some (arbitrarily chosen) vertex of $`G`$. This can be accomplished by means of a simple Horn program. Let us choose a vertex, say $`s`$, in $`G`$. Include in $`T_H`$ the Horn rules $`hc(s,t)vstd(t)`$, for every edge $`(s,t)`$ in $`G`$. In addition, include in $`T_H`$ Horn rules $`vstd(t),hc(t,u)vstd(u)`$, for every edge $`(t,u)`$ of $`G`$ not starting in $`s`$. Clearly, the least model of $`FT_H`$, where $`F`$ is a subset of $`E`$, contains precisely these variables of the form $`vstd(t)`$ for which $`t`$ is reachable from $`s`$ by a nonempty path spanned by the edges in $`F`$. Thus, $`F`$ is the set of edges of a hamilton cycle of $`G`$ if and only if the least model of $`FT_H`$, contains variable $`vstd(t)`$ for every vertex $`t`$ of $`G`$. Let us define $`T_{PC}=\{vstd(t):tV\}`$ and $`T_{ham}(G)=(T_C,T_H,T_{PC})`$. It follows that hamilton cycles of $`G`$ can be reconstructed (in linear time) from answer sets to the $`\mathrm{DC}`$ theory $`T_{ham}(G)`$. In other words, to find a hamilton cycle in $`G`$, it is enough to find an answer set for $`T_{ham}(G)`$.

This example illustrates the simplicity of the semantics — it is only a slight adaptation of the semantics of propositional logic to the case when in addition to propositional clauses we also have Horn rules in theories. It also illustrates the power of $`\mathrm{DC}`$ to generate concise encodings. All known propositional encodings of the hamilton-cycle problem require that additional variables are introduced to “count” how far from the starting vertex an edge is located. Consequently, propositional encodings are much larger and lead to inefficient computational approaches to the problem. We present experimental evidence to this claim later in the paper.

The question arises which search problems can be represented (and solved) by means of finding answer sets to appropriate $`\mathrm{DC}`$ programs. In general, the question remains open. We have an answer, though, if we restrict our attention to the special case of decision problems. Consider a $`\mathrm{DC}`$ theory $`T=(T_C,T_H,T_{PC})`$, where $`T_H=T_{PC}=\mathrm{}`$. Clearly, $`M`$ is an answer set for $`T`$ if and only if $`M`$ is a model of the collection of clauses $`T_C`$. Thus, the problem of existence of an answer set is at least as hard as the propositional satisfiability problem. On the other hand, for every $`\mathrm{DC}`$ theory $`T`$ and for every set $`MAt_C(T)`$, it can be checked in linear time whether $`M`$ is an answer set for $`T`$. Thus, we obtain the following complexity result.

###### Theorem 1

The problem of existence of an answer set for a finite propositional $`\mathrm{DC}`$ theory $`T`$ is NP-complete.

It follows that every problem in NP can be polynomially reduced to the problem of existence of an answer set for a propositional $`\mathrm{DC}`$ program. Thus, given a problem $`\mathrm{\Pi }`$ in NP, for every instance $`I`$ of $`\mathrm{\Pi }`$, $`\mathrm{\Pi }`$ can be decided by deciding the existence of an answer set for the $`\mathrm{DC}`$ program corresponding to $`\mathrm{\Pi }`$ and $`I`$.

Propositional $`\mathrm{DC}`$ can be extended to the predicate case. It is important as it significantly simplifies the task of developing programs for solving problems with $`\mathrm{DC}`$. In the example discussed above, the theory $`T_{ham}(G)`$ depends heavily on the input. Each time we change the input graph, a different theory has to be used. However, when constructing predicate $`\mathrm{DC}`$-based solutions to a problem $`\mathrm{\Pi }`$, it is often possible to separate the representation of an instance (input) to $`\mathrm{\Pi }`$ from that of the constraints that define $`\mathrm{\Pi }`$. As a result only one (predicate) program describing the constraints of $`\mathrm{\Pi }`$ needs to be written. Specific input for the program, say $`I`$, can be described separately as a collection of facts (according to some uniform schema). Both parts together can be combined to yield a $`\mathrm{DC}`$ program whose answer sets determine solutions to $`\mathrm{\Pi }`$ for the input $`I`$. Such an approach, we will refer to it as uniform, is often used in the context of DATALOG, DATALOG<sup>¬</sup> or logic programming to study complexity of these systems as query languages. The part representing input is referred to as the extensional database. The part representing the query or the problem is called the intensional database or program. Due to the space limitations we do not discuss the details of the predicate case here. They will be given in the full version of the paper. We only state a generalization of Theorem 1.

###### Theorem 2

The expressive power of $`\mathrm{DC}`$ is the same as that of stable logic programming. In particular, a decision problem $`\mathrm{\Pi }`$ can be solved uniformly in $`\mathrm{DC}`$ if and only if $`\mathrm{\Pi }`$ is in the class NP.

## Implementation

Some types of constraints appear frequently in applications. For instance, when defining plans we may want to specify a constraint that says that exactly one action from the set of allowed actions be selected at each step. Such constraints can be modeled by collections of clauses. To make sure $`\mathrm{DC}`$ programs are as easy to write and as concise as possible we have extended the syntax of $`\mathrm{DC}`$ by providing explicit ways to model constraints of the form “select at least (at most, exactly) $`k`$ elements from a set”. Having these constraints results in shorter programs which, as we believe, has a significant positive effect of the performance of our system.

An example of a select constraint with a short explanation is presented here. Let $`PRED`$ be the set of predicates occurring in the IDB. For each variable $`X`$ declared in the IDB the range $`R(X)`$ of $`X`$ is determined by the EDB.

where $`n,m`$ are nonnegative integers such that $`nm,qPRED`$ and $`p_1,\mathrm{},p_i`$ are EDB predicates or logical conditions (logical conditions can be comparisons of arithmetic expressions or string comparisons). The interpretation of this constraint is as follows: for every $`\stackrel{}{x}R(\stackrel{}{X})`$ at least $`n`$ atoms and at most $`m`$ atoms in the set $`\{q(\stackrel{}{x},\stackrel{}{y}):\stackrel{}{y}R(\stackrel{}{Y})\}`$ are true.

We implemented $`\mathrm{DC}`$ in the predicate setting. Thus, our system consists of two main modules. The first of them, referred to as grounder, converts a predicate $`\mathrm{DC}`$ program (consisting of both the extensional and intensional parts) into the corresponding propositional $`\mathrm{DC}`$ program. The second module, $`\mathrm{DC}`$ solver, denoted dcs, finds the answer sets to propositional $`\mathrm{DC}`$ programs. Since we focus on the propositional case here, we only describe the key ideas behind the $`\mathrm{DC}`$ solver, dcs.

The $`\mathrm{DC}`$ solver uses a Davis-Putnam type approach, with backtracking, propagation and lookahead (also called literal testing), to deal with constraints represented as clauses, select constraints and Horn rules, and to search for answer sets. The lookahead in $`\mathrm{DC}`$ is similar to local processing performed in csat (?). However, we use different methods to determine how many literals to consider in the lookahead phase. Other techniques, especially propagation and search heuristics, were designed specifically for the case of $`\mathrm{DC}`$ as they must take into account the presence of Horn rules in programs.

The lookahead procedure selects a number of literals which have not yet been assigned a value. For each such literal, the procedure tries both truth values: true and false. For each assignment, the theory is evaluated using propagation. If in both cases a contradiction is reached, then it is necessary to backtrack. If for only one evalution a conflict is reached, then the literal is assigned the other truth value and we proceed to the next step. If neither evaluation results in a contradiction, we cannot assign a truth value to this literal but we save the data such as the number of forced literals and the number of clauses satisfied, computed during propagation.

Clearly, if all unassigned literals were tested it would prune the most search space. At the same time, the savings might not be large enough to compensate for the increase in the running time caused by extensive lookahead. Thus, we select only a portion of all unassigned literals for lookahead. The number of literals to consider was established empirically (it does not depend on the size of the theory). Since not all literals are selected, it is important to focus on those literals that are likely to result in a contradiction for at least one of the truth values. In our implementation, we select the most constrained literals, as determined by their weights.

Specifically, each constraint is assigned a weight based on its current length and types (recall that in addition to propositional clauses, we also allow other types of constraints, e.g., select constraints). The shorter the constraint the greater its weight. Also, certain types of constraints force more assignments on literals and are given a greater weight than other constraints of the same length. Every time a literal appears in an unsatisfied constraint, the weight of that literal is incremented by the weight of the clause.

After testing a predetermined number of literals without finding a forced truth assignment and without backtracking, the information computed during propagation is used to choose the next literal for which both possible truth assignments have to be tested (branching literal). The choice of the next branching literal is based on an approximation of which literal, once assigned a truth value, will force the truth assignments onto the largest number of other literals and will satisfy the largest number of clauses. Using the data computed during propagation gives more accurate information on which to base such approximations. The methods used for determining which literals to select in the lookahead and which data to collect and save during the propagation phase are two key ways in which the literal testing procedure differs from the local processing of csat.

## Experimentation

We compared the performance of $`\mathrm{DC}`$ solver dcs with smodels, a system for computing stable models of logic programs (?), and csat, a system for testing propositional satisfiability (?). In the case of smodels we used version 2.24 in conjunction with the grounder lparse, version 0.99.41. These versions of lparse and smodels implement the expressive rules described by (?). The expressive rules were used whenever applicable during the testing. The programs were all executed on a Sun SparcStation 20. For each test we report the cpu user times for processing the corresponding propositional program or theory. We tested all three system to compute hamilton cycles and colorings in graphs, to solve the $`N`$-queens problem, to prove that the pigeonhole problem has no solution if the number of pigeons exceeds the number of holes, and to compute Schur numbers.

The Hamilton cycle problem has already been described. We randomly generated one thousand graphs with the edge-to-vertex ratio such that $`50\%`$ of the graphs contained Hamilton cycles (crossover region). The number of vertices ranged from 30 to 80. We used encodings of the problem as a $`\mathrm{DC}`$ program, logic program (in smodels syntax) and as a propositional theory. dcs performed better than smodels and smodels performed significantly better than csat (Fig. 1). We believe that a major factor behind poorer performance of csat is that all known propositional encodings of the hamilton cycle problem are much larger than those possible with $`\mathrm{DC}`$ or logic programs (under the stable model semantics). Propositional encodings, due to their size, rendered csat not practical to execute for graphs with more than 40 vertices.

The $`N`$-queens problem consists of finding a placement of $`N`$ queens on an $`N\times N`$ board such that no queen can remove another. Both csat and dcs execute in much less time than smodels for these problems (Fig. 2). Again the size of the encoding seems to be a major factor. One thing to consider in this case is that the number of rules for smodels is approximately five times that for $`\mathrm{DC}`$ and more than twice that of propositional encodings.

The Schur problem consists of placing $`N`$ numbers $`1,2,\mathrm{},N`$ in $`B`$ bins such that no bin is closed under sums. That is, for all numbers $`x`$, $`y`$, $`z`$, $`1x,y,zN`$, if $`x`$ and $`y`$ are the same bin, then $`z`$ is not ($`x`$ and $`y`$ need not be distinct). The Schur number $`S(B)`$ is the maximum number $`N`$ for which such a placement is still possible. It is known to exist for every $`B1`$. We considered the problem of the existence of the placement for $`B=3`$ and $`N=13`$ and $`14`$, and for $`B=4`$ and $`N=43,44`$ and $`45`$. In each case we used all three systems to process the corresponding encodings. The results are shown in Fig. 3. It follows that $`S(3)=13`$ and $`S(4)=44`$. Again, dcs outperforms both smodels and csat.

Results for graph 3-coloring for graphs with the number of vertices ranging from $`50`$ to $`300`$ are shown in Fig. 4 (for every choice of the number of vertices, 100 graphs from the crossover region were randomly generated). Both dcs and csat performed better than smodels. Again the size of the theory seems to be a factor. The CNF theory for coloring is smaller than a logic program encoding the same problem. The sizes of propositional and $`\mathrm{DC}`$ encodings are similar.

Results for the pigeonhole placement problem show a similar performance of all three algorithms, with csat doing slightly better than the others and dcs outperforming (again only slightly) smodels.

## Conclusions

We described a new system, $`\mathrm{DC}`$, for solving search problems. We designed $`\mathrm{DC}`$ so that its semantics was as close as possible to that of propositional logic. Our goal was to design a system that would result is short problem encodings. Thus, we provided constructs for some frequently occurring types of constraints and we built into $`\mathrm{DC}`$ elements of nonmonotonicity by including Horn rules in the syntax. As a result, $`\mathrm{DC}`$ programs encoding search problems are often much smaller than those possible with propositional theories. Experimental results show that dcs often outperforms systems based on propositional satisfiability as well as systems based on nonmonotonic logics, and that it constitutes a viable approach to solving problems in AI, constraint satisfaction and combinatorial optimization. We believe that our focus on short programs is the key to the success of $`\mathrm{DC}`$ and its reasoning engine dcs. Our results show that when building general purpose solvers of search problems, the size of encodings should be a key design factor.