format: text

text:
# Untitled Document

Rectangular Scott-type Permanents

Guo-Niu Han et Christian Krattenthaler

Abstract.— Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of a polynomial $`P(x)`$ of degree $`n`$ and $`y_1,y_2,\mathrm{},y_m`$ be the zeroes of another polynomial $`Q(y)`$ of degree $`m`$. Our object of study is the permanent $`\text{per}(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$, here named “Scott-type” permanent, the case of $`P(x)=x^n1`$ and $`Q(y)=y^n+1`$ having been considered by R. F. Scott. We present an efficient approach to determining explicit evaluations of Scott-type permanents, based on generalizations of classical theorems by Cauchy and Borchardt, and of a recent theorem by Lascoux. This continues and extends the work initiated by the first author (“Généralisation de l’identité de Scott sur les permanents,” to appear in Linear Algebra Appl.). Our approach enables us to provide numerous closed form evaluations of Scott-type permanents for special choices of the polynomials $`P(x)`$ and $`Q(y)`$, including generalizations of all the results from the above mentioned paper and of Scott’s permanent itself. For example, we prove that if $`P(x)=x^n1`$ and $`Q(y)=y^{2n}+y^n+1`$ then the corresponding Scott-type permanent is equal to $`(1)^{n+1}n!`$.

Résumé.—Soient $`x_1,x_2,\mathrm{},x_n`$ les zéros d’un polynôme $`P(x)`$ de degré $`n`$ et $`y_1,y_2,\mathrm{},y_m`$ les zéros d’un autre polynôme $`Q(y)`$ de degré $`m`$. Notre objet d’étude est le permanent $`\text{per}(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$, appelé ici permanent de type Scott. Le cas de $`P(x)=x^n1`$ et $`Q(y)=y^n+1`$ a été considéré par R. F. Scott. Nous présentons une approche efficace pour déterminer les évaluations explicites des permanents de type Scott, basée sur des généralisations des théorèmes classiques de Cauchy et Borchardt, et d’un théorème récent de Lascoux. La présente étude prolonge le travail du premier auteur (“Généralisation de l’identité de Scott sur les permanents,” à apparaître dans Linear Algebra Appl.). Notre approche nous permet de fournir de nombreuses évaluations explicites des permanents de type Scott pour des choix spéciaux des polynômes $`P(x)`$ et $`Q(y)`$, y compris des généalisations de tous les résultats de l’article mentionné ci-dessus et du permanent de Scott lui-même. Par exemple, nous prouvons que si $`P(x)=x^n1`$ et $`Q(y)=y^{2n}+y^n+1`$ alors le permanent correspondant de type Scott est égal à $`(1)^{n+1}n!`$.

1. Introduction

In 1881, Scott  stated, without proof, the following result:

Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$ and $`y_1,y_2,\mathrm{},y_n`$ be the zeroes of $`y^n+1`$. Let $`A`$ be the $`n\times n`$ matrix $`(1/(x_iy_j))_{1i,jn}`$. Then

$$\text{per}(A)=\{\begin{array}{cc}(1)^{\frac{n1}{2}}\frac{n(135\mathrm{}(n2))^2}{2^n},\hfill & \text{if }n\text{ is odd,}\hfill \\ 0,\hfill & \text{if }n\text{ is even.}\hfill \end{array}$$

In 1978, in his monograph Permanents, Minc \[13, p. 155\] included this result in a list of conjectures on permanents. Since then, several proofs have been given , one of which by Minc himself. All of these proofs are heavily based on the fact that the zeroes of the polynomials $`x^n1`$ and $`y^n+1`$ can be written in simple explicit terms. Thus, neither of these proofs extends to the more general problem which is the subject of this paper:

Given a polynomial $`P(x)`$ of degree $`n`$ and a polynomial $`Q(y)`$ of degree $`m`$, let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`P(x)`$ and $`y_1,y_2,\mathrm{},y_m`$ be the zeroes of $`Q(y)`$. Evaluate the permanent of the $`n\times m`$ matrix $`(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$.

As usual, the permanent $`\text{per}(A)`$ of an $`n\times m`$ matrix $`A`$ is defined as the sum of all possible products of $`n`$ coefficients of $`A`$ chosen such that no two of the coefficients are taken from the same row nor from the same column (see \[13, Ch. 1, (1.1)\]). Given the assumptions of the above problem, we call the permanent of the matrix $`(1/(x_iy_j))_{1in,\mathrm{\hspace{0.17em}1}jm}`$ a Scott-type permanent, and denote it by $`\text{PER}(P(x),Q(y))`$.

In , the first author presented a new approach to this type of problem in the case $`n=m`$, i.e., in the case that both polynomials have the same degree. This approach does not rely at all on explicit analytic forms of zeroes of polynomials. Instead, it makes essential use of recent symmetric functions techniques, in particular of a theorem due to Lascoux , which the latter author established in his étude on the square ice model of statistical mechanics.

In the present paper, we are going to extend this approach to arbitrary $`n`$ and $`m`$. This requires extensions of classical theorems of Cauchy and Borchardt (see Theorems (Cauchy+) and (Borchardt+) in Section 2), and an extension of Lascoux’s theorem (see Theorem (Lascoux+)). As a result (see Theorem 1), we are able to express any Scott-type permanent as the quotient of a determinant which features complete homogeneous and elementary symmetric functions in the zeroes of the two polynomials, divided by the resultant of the two polynomials. In particular, it follows immediately that any Scott-type permanent is rational in the coefficients of the polynomials $`P(x)`$ and $`Q(y)`$.

In Section 5 we apply this result to obtain explicit evaluations of Scott-type permanents in numerous special cases. Amongst others, we provide generalizations of all the results from , thus also covering Scott’s permanent itself. For the proofs of the results in Section 5, we make use of two particular specializations of our main theorem, Theorem 1, which we derive in Section 3 (see Theorems 3 and 4), and of four determinant evaluations, which we state and establish separately in Section 4.

Finally, we also comment briefly on an alternative approach to the evaluation of Scott-type determinants, due to the first author . It allows to express Scott-type permanents in terms of weighted sums over involutions (see Theorem 2 in Section 2). By combining this result with some of the evaluations in Section 5, we obtain interesting summation theorems, which are presented in Section 6.

2. The general theory

In , the main ingredients are theorems by Cauchy, Borchardt, Lascoux, and a lemma on the resultant. Since we intend to extend the approach of  to the case of rectangular Scott-type permanents (corresponding to polynomials of, possibly, different degrees), we have to first provide the appropriate extensions of these theorems.

Given positive integers $`m`$ and $`n`$ and two sets $`X=\{x_1,x_2,\mathrm{},x_n\}`$ and $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ of variables, we use the following notations:

$$R(X,Y):=\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{m}{}}(x_iy_j)\text{ and }\mathrm{\Delta }(X):=\underset{i<j}{}(x_ix_j).$$

$$\left(\frac{1}{x_iy_j}\right):=\left(\begin{array}{cccc}\frac{1}{x_1y_1}& \frac{1}{x_1y_2}& \mathrm{}& \frac{1}{x_1y_m}\\ \frac{1}{x_2y_1}& \frac{1}{x_2y_2}& \mathrm{}& \frac{1}{x_2y_m}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \frac{1}{x_ny_1}& \frac{1}{x_ny_2}& \mathrm{}& \frac{1}{x_ny_m}\end{array}\right).$$

We first state a variation on Cauchy’s evaluation of his double alternant (cf. \[4; 15, vol. 1, pp. 342–345\]).

Theorem (Cauchy+).— For $`mn`$, let

$$C(X,Y):=\left(\begin{array}{cccc}\frac{1}{x_1y_1}& \frac{1}{x_1y_2}& \mathrm{}& \frac{1}{x_1y_m}\\ \frac{1}{x_2y_1}& \frac{1}{x_2y_2}& \mathrm{}& \frac{1}{x_2y_m}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \frac{1}{x_ny_1}& \frac{1}{x_ny_2}& \mathrm{}& \frac{1}{x_ny_m}\\ \\ 1& 1& \mathrm{}& 1\\ y_1& y_2& \mathrm{}& y_m\\ y_1^2& y_2^2& \mathrm{}& y_m^2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ y_1^{mn1}& y_2^{mn1}& \mathrm{}& y_m^{mn1}\end{array}\right).$$

Then

$$det(C(X,Y))=(1)^{n(n1)/2}\frac{\mathrm{\Delta }(X)\mathrm{\Delta }(Y)}{R(X,Y)}.$$

Proof.—If $`n=m`$, this is exactly Cauchy’s theorem. The general case can be either established directly, or, it may be observed that the “general” case is in fact implied by Cauchy’s theorem. To see this, consider the above identity with $`n=m`$. Given $`k<m`$, expand both sides as power series in $`1/x_{k+1}`$, …, $`1/x_m`$, and compare coefficients of $`1/x_{k+1}x_{k+2}^2\mathrm{}x_m^{mk}`$ on both sides.

Next we state the required extension of Borchardt’s theorem \[2; 15, vol. 2, pp. 173–175\]. It can be established by reading through the proof of Borchardt’s theorem given in \[1, Proof of Cor. 5.1\], ignoring however the restriction $`m=n`$ (see also ).

Theorem (Borchardt+).— For $`mn`$, let

$$B(X,Y):=\left(\begin{array}{cccc}\frac{1}{(x_1y_1)^2}& \frac{1}{(x_1y_2)^2}& \mathrm{}& \frac{1}{(x_1y_m)^2}\\ \frac{1}{(x_2y_1)^2}& \frac{1}{(x_2y_2)^2}& \mathrm{}& \frac{1}{(x_2y_m)^2}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \frac{1}{(x_ny_1)^2}& \frac{1}{(x_ny_2)^2}& \mathrm{}& \frac{1}{(x_ny_m)^2}\\ \\ 1& 1& \mathrm{}& 1\\ y_1& y_2& \mathrm{}& y_m\\ y_1^2& y_2^2& \mathrm{}& y_m^2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ y_1^{mn1}& y_2^{mn1}& \mathrm{}& y_m^{mn1}\end{array}\right).$$

Then

$$det(B(X,Y))=det(C(X,Y))\times \text{per}\left(\frac{1}{x_iy_j}\right).$$

Our next goal is to derive the required extension of (a special case of) Lascoux’s theorem \[11, Theorem q\]. Let $`Z=\{z_1,z_2,\mathrm{},z_n\}`$ be another set of variables, of equal cardinality as $`X`$. For $`1in`$ we define the divided difference $`_i`$ by

$$_i:f\frac{ff^{\sigma _i}}{x_iz_i},$$

where $`\sigma _i`$ is the transposition which interchanges $`x_i`$ and $`z_i`$. It is easy to see that

$$_i\frac{1}{z_iy}=\frac{1}{(x_iy)(z_iy)}.$$

Since the operator $`_i`$ acts only on one row of the matrix $`C(Z,Y)`$ (to be precise, the $`i`$-th row), it follows that

$$_1_2\mathrm{}_n(det(C(Z,Y))|_{Z=X}=det(B(X,Y)).$$
$`(1)`$
Now, to generalize Lascoux’s theorem to our case, one simply readsthrough the proof of Theorem q in , on introducing slight modifications if necessary. The result is:

Theorem (Lascoux+).— Let $`H(X)`$ be the $`n\times (m+n1)`$ matrix defined by

$$H(X):=\left(h_{ji}(X)\right)_{1in,1jm+n1},$$

where $`h_s(X)`$ denotes the complete homogeneous symmetric function of degree $`s`$ in the variables $`X`$ (cf. \[12, Ch. 1\]), and $`E(Y)`$ be the $`(m+n1)\times n`$ matrix defined by

$$E(Y):=\left((j2k+2)(1)^{mj+k1}e_{mj+k1}(Y)\right)_{1jm+n1,1kn},$$

where $`e_s(Y)`$ denotes the elementary symmetric function of degree $`s`$ in the variables $`Y`$ (cf. \[12, Ch. 1\]). Then

$$_1_2\mathrm{}_n(\mathrm{\Delta }(Z)R(X,Y))|_{Z=X}=\mathrm{\Delta }(X)det\left(H(X)E(Y)\right).$$

Now we are in the position to state our main theorem, which will enable us, in Section 5, to evaluate numerous Scott-type permanents in closed form. The theorem implies immediately that any Scott-type permanent $`\text{PER}(P(x),Q(y))`$ is a rational function in the coefficients of the two polynomials $`P(x)`$ and $`Q(y)`$.

Theorem 1.— Let $`m`$ and $`n`$ be arbitrary positive integers, and let $`X=\{x_1,x_2,\mathrm{},x_n\}`$ and $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ be two sets of variables. Then

$$\text{per}\left(\frac{1}{x_iy_j}\right)=\frac{det\left(H(X)E(Y)\right)}{R(X,Y)}.$$

Proof.—First let $`mn`$. By combining Theorems (Cauchy+) and (Borchardt+), and Equation (1), we obtain

$$(1)^{n(n1)/2}\frac{R(X,Y)}{\mathrm{\Delta }(X)\mathrm{\Delta }(Y)}_1_2\mathrm{}_n(det(C(Z,Y))|_{Z=X}$$

for the permanent. Now we apply Theorem (Cauchy+) again in order to replace $`C(Z,Y)`$ by the corresponding product form guaranteed by the theorem. After having used that the divided differences $`_i`$ commute with $`\mathrm{\Delta }(Y)/R(X,Y)R(Z,Y)`$ (because the latter expression is symmetric in $`x_i`$ and $`z_i`$), Theorem (Lascoux+) applies and yields the desired result.

If $`m<n`$, the permanent clearly vanishes. According to Theorem (Lascoux+), it suffices to establish that

$$U:=_1_2\mathrm{}_n(\mathrm{\Delta }(Z)R(X,Y))=0.$$

To begin with, we rewrite the Vandermonde determinant evaluation as $`\mathrm{\Delta }(Z)R(X,Y)=(1)^{n(n1)/2}det(z_i^{j1}R(x_i,Y))`$. Writing $`Q(x_i):=R(x_i,Y)=a_mx_i^m+\mathrm{}+a_1x_i+a_0`$, we obtain for $`U`$ the expression

$$U=det(_iz_i^{j1}R(x_i,Y))=det\left(z_i^{j1}Q(x_i)x_i^{j1}Q(z_i)\right)/(x_iz_i).$$

Since for $`m<n`$, we have

$$\underset{j=0}{\overset{m}{}}a_m\left(z_i^jQ(x_i)x_i^jQ(z_i)\right)=0,$$

the $`m+1`$ elements $`z_i^jQ(x_i)x_i^jQ(z_i)`$, $`0jm`$, are linearly dependent. Hence, $`U=0`$.

In , the first author obtained another expression for the permanent, in form of a certain weighted sum over involutions. To state and explain this formula, for $`sX`$ define

$$L(s;X,Y):=\underset{xs}{}\frac{1}{xs}+\underset{yY}{}\frac{1}{sy}.$$

Let us denote by $`(n)`$ the set of involutions on $`\{1,2,\mathrm{},n\}`$. Given an involution $`\sigma (n)`$, we define the weight $`\mathrm{\Psi }(\sigma )`$ of $`\sigma `$ by

$$\mathrm{\Psi }(\sigma ;X,Y):=\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}L(x_k;X,Y)$$

where the first product is over all transpositions $`(ij)`$ in the disjoint cycle decomposition of $`\sigma `$, and where the second product is over all fixed points $`k`$ of $`\sigma `$. Then the result from  is the following.

Theorem 2.— Let $`m`$ and $`n`$ be arbitrary positive integers, and let $`X=\{x_1,x_2,\mathrm{},x_n\}`$ and $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ be two sets of variables. Then

$$\text{per}\left(\frac{1}{x_iy_j}\right)=\underset{\sigma (n)}{}\mathrm{\Psi }(\sigma ;X,Y).$$

Example 1.— Let $`n=1`$, $`X=\{x\}`$, $`Y=\{y_1,y_2,\mathrm{},y_m\}`$. Then

$$\begin{array}{cc}\hfill H(X)& :=\left(h_{j1}(X)\right)_{i=1,1jm}=(h_0(x),h_1(x),\mathrm{},h_{m1}(x)),\hfill \\ \hfill E(Y)& :=\left(j(1)^{mj}e_{mj}(Y)\right)_{1jm,k=1}\hfill \\ & =(1(1)^{m1}e_{m1}(Y),2(1)^{m2}e_{m2}(Y),\mathrm{},me_0(Y))^t.\hfill \end{array}$$

We have

$$H(X)E(Y)=\underset{j=0}{\overset{m1}{}}(j+1)(1)^{mj1}e_{mj1}(Y)x^j.$$

Therefore,

$$\begin{array}{cc}\hfill \text{per}\left(\frac{1}{xy}\right)& =\frac{1}{xy_1}+\frac{1}{xy_2}+\mathrm{}+\frac{1}{xy_m}\text{[Definition, Th. 2]}\hfill \\ & =\frac{\underset{j=0}{\overset{m1}{}}(j+1)(1)^{mj1}e_{mj1}(Y)x^j}{(xy_1)(xy_2)\mathrm{}(xy_m)}.\text{[Th. 1]}\hfill \end{array}$$

Example 2.— For $`n=4`$ and $`m=3`$, Theorem 1 yields the following identity:

$$det\left(\left(\begin{array}{cccccc}h_0& h_1\hfill & h_2\hfill & h_3\hfill & h_4\hfill & h_5\hfill \\ 0& h_0\hfill & h_1\hfill & h_2\hfill & h_3\hfill & h_4\hfill \\ 0& 0\hfill & h_0\hfill & h_1\hfill & h_2\hfill & h_3\hfill \\ 0& 0\hfill & 0\hfill & h_0\hfill & h_1\hfill & h_2\hfill \end{array}\right)\times \left(\begin{array}{cccc}e_2& e_3\hfill & 0\hfill & 0\hfill \\ 2e_1& 0\hfill & 2e_3\hfill & 0\hfill \\ 3e_0& e_1\hfill & e_2\hfill & 3e_3\hfill \\ 0& 2e_0\hfill & 0\hfill & 2e_2\hfill \\ 0& 0\hfill & e_0\hfill & e_1\hfill \\ 0& 0\hfill & 0\hfill & 0\hfill \end{array}\right)\right)=0.$$

For $`n=2`$ and $`m=1`$, Theorem 2 yields

$$\frac{1}{(x_1x_2)^2}+\left(\frac{1}{x_2x_1}+\frac{1}{x_1y}\right)\left(\frac{1}{x_1x_2}+\frac{1}{x_2y}\right)=0.$$

3. The case of $`P(x)=x^n1`$ and of $`P(x)=x^{n1}+\mathrm{}+x+1`$

In this section we specialize Theorem 1 to the case that the $`x_i`$’s are the zeroes of the polynomial $`P(x)=x^n1`$ or of $`P(x)+x^{n1}+\mathrm{}+x+1`$, and the $`y_i`$’s are the zeroes of an arbitrary other polynomial. (This covers, for example, the case of Scott’s identity). For the remainder of this section, we fix $`m`$ and $`n`$, $`mn`$.

Let $`X=\{x_1,x_2,\mathrm{},x_n\}`$ be the set of zeroes of $`x^n1`$, and let $`Y=\{y_1,y_2,\mathrm{},y_m\}`$ be the set of zeroes of $`Q(x)=a_my^m+a_{m1}y^{m1}+\mathrm{}+a_1y^1+a_0`$, with $`a_m=1`$. We write

$$\text{PER}(P,Q):=\text{per}\left(\frac{1}{x_iy_j}\right)=\frac{det(H(X)E(Y))}{R(X,Y)}.$$
$`(2)`$
Since

$$h_i(X)t^i=\frac{1}{_i(1tx_i)}=\frac{1}{(1t^n)}=1+t^n+t^{2n}+\mathrm{},$$

we have

$$h_k(X)=\{\begin{array}{cc}1,\hfill & \text{if }k=0(\text{mod }n)\text{,}\hfill \\ 0,\hfill & \text{if }k0(\text{mod }n)\text{.}\hfill \end{array}$$

We denote by $`\text{I}_k`$ the $`k\times k`$ identity matrix, and by $`\text{0}_{l,c}`$ the $`l\times c`$ matrix with all entries equal to $`0`$. For all $`r`$, we write $`r\%n`$ for the number between $`1`$ and $`n`$ that satisfies $`r(\text{mod }n)=r\%n(\text{mod }n)`$. Then we have

$$H(X)=\left(\begin{array}{ccccc}\text{I}_n& \text{I}_n& \mathrm{}& \text{I}_n& \begin{array}{c}\text{I}_m^{}\\ \text{0}_{nm^{},m^{}}\end{array}\end{array}\right),$$

with $`m^{}=(m+n1)\%n`$.

Furthermore, let $`\mathrm{diag}_n^i(c_1,c_2,\mathrm{},c_n)`$ denote the $`n\times n`$ “diagonal” matrix, in which the (broken) diagonal starts in the $`i`$-th row,

$$\left(\begin{array}{cc}\text{0}_{i1,ni+1}& \mathrm{diag}_{i1}(c_{ni+2},c_{ni+3},\mathrm{},c_n)\\ \mathrm{diag}_{ni+1}(c_1,c_2,\mathrm{},c_{ni+1})& \text{0}_{ni+1,i1}\end{array}\right).$$

According to the definition of $`E(Y)`$, a simple calculation yields that

$$H(X)E(Y)=\underset{r=0}{\overset{m}{}}\mathrm{diag}_n^{r\%n}(ra_r,(r1)a_r,\mathrm{},(rn+1)a_r).$$

Example 3.— For $`n=3`$ and $`m=4`$, let $`P(x)=x^31`$ and $`Q(y)=y^4+a_3y^3+a_2y^2+a_1y+a_0`$. Then $`H(X)E(Y)`$ is the sum of the following matrices:

$$\left(\begin{array}{ccc}0& a_0\hfill & 0\hfill \\ 0& 0\hfill & 2a_0\hfill \\ 0& 0\hfill & 0\hfill \end{array}\right)+\left(\begin{array}{ccc}a_1& 0\hfill & 0\hfill \\ 0& 0\hfill & 0\hfill \\ 0& 0\hfill & a_1\hfill \end{array}\right)+\left(\begin{array}{ccc}0& 0\hfill & 0\hfill \\ 2a_2& 0\hfill & 0\hfill \\ 0& a_2\hfill & 0\hfill \end{array}\right)+\left(\begin{array}{ccc}0& 2a_3\hfill & 0\hfill \\ 0& 0\hfill & a_3\hfill \\ 3a_3& 0\hfill & 0\hfill \end{array}\right)+\left(\begin{array}{ccc}4& 0\hfill & 0\hfill \\ 0& 3\hfill & 0\hfill \\ 0& 0\hfill & 2\hfill \end{array}\right).$$

In the theorem below, we summarize our findings.

Theorem 3.— Let $`P(x)=x^n1`$ and $`Q(y)=a_my^m+\mathrm{}+a_1y^1+a_0`$, $`a_m`$ not necessarily $`1`$. Writing

$$\text{Fes}(Q)=det\left(\underset{r=0}{\overset{m}{}}\mathrm{diag}_n^{r\%n}(ra_r,(r1)a_r,\mathrm{},(rn+1)a_r)\right),$$
$`(3)`$
we have

$$\text{PER}(P,Q)=\frac{\text{Fes}(Q)}{\text{Res}(P,Q)},$$

where Res is the classical resultant of two polynomials.

Proof.—We have already seen that the theorem is true for $`a_m=1`$. On the other hand, there hold $`\text{Fes}(\lambda Q)=\lambda ^n\text{Fes}(Q)`$ and $`\text{Res}(P,\lambda Q)=\lambda ^n\text{Res}(P,Q)`$, as is easily verified.

Now let us consider the case that $`X=\{x_1,x_2,\mathrm{},x_{n1}\}`$ is the set of zeroes of $`P(x)=x^{n1}+\mathrm{}+x+1`$. We perform an analysis very similar to the one before, using the fact that we have

$$h_i(X)t^i=\frac{1}{_i(1tx_i)}=\frac{1t}{_i(1t^n)}=1t+t^nt^{n+1}+\mathrm{}.$$

In order to state the result, we introduce the following notation: We write $`\stackrel{𝑒}{\mathrm{diag}}{}_{n1}{}^{i}(c_1,c_2,\mathrm{},c_{n1})`$ for the $`(n1)\times (n1)`$ matrix

$$\left(\begin{array}{ccc}\text{0}_{i1,ni}& \text{0}_{i1,1}& \mathrm{diag}_{i2}(c_{ni+2},c_{ni+3},\mathrm{},c_{n1})\\ \mathrm{diag}_{ni}(c_1,c_2,\mathrm{},c_{ni})& \text{0}_{ni,1}& \text{0}_{ni+1,i2}\end{array}\right)$$

if $`i>1`$, respectively $`\mathrm{diag}_{n1}(c_1,c_2,\mathrm{},c_{n1})`$ if $`i=1`$. This is again a matrix with a (possibly broken) diagonal, in which the diagonal “jumps over” one row and column in the case that it is broken. (Note the slight discrepancy in dimension between the diagonal and the zero matrices.)

Theorem 4.— Let $`P(x)=x^{n1}+\mathrm{}+x+1`$ and $`Q(y)=a_my^m+\mathrm{}+a_1y^1+a_0`$, $`a_m`$ not necessarily $`1`$. Writing

$$\begin{array}{cc}\hfill \stackrel{𝑒}{\mathrm{Fes}}(Q)=& det(\underset{r=0}{\overset{m}{}}\stackrel{𝑒}{\mathrm{diag}}{}_{n1}{}^{r\%n}(ra_r,(r1)a_r,\mathrm{},(rn+2)a_r)\hfill \\ & \underset{r=0}{\overset{m}{}}\stackrel{𝑒}{\mathrm{diag}}{}_{n1}{}^{(r1)\%n}(ra_r,(r1)a_r,\mathrm{},(rn+2)a_r)),\hfill \end{array}$$
$`(4)`$
we have

$$\text{PER}(P,Q)=\frac{\stackrel{𝑒}{\mathrm{Fes}}(Q)}{\text{Res}(P,Q)},$$

where, again, Res is the classical resultant of two polynomials.

According to Theorems 3 and 4, for accomplishing the evaluation of the permanent, it is necessary to evaluate the numerator $`\text{Fes}(Q)`$, respectively $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, and the denominator $`\text{Res}(P,Q)`$. For the evaluation of the resultant $`\text{Res}(P,Q)`$, we make use of the following lemma, which, for example, appears explicitly in . (In fact, it follows from a special case of Proposition 6 in the next section.)

Lemma 5.—Let $`d`$ be the greatest common divisor of $`m`$ and $`n`$, and let $`A`$ and $`C`$ be two nonzero constants. Then the resultant of the two polynomials $`Ax^mB`$ and $`Cx^nD`$ is given by

$$\text{Res}(Ax^mB,Cx^nD)=(1)^m\left(A^{n/d}D^{m/d}B^{n/d}C^{m/d}\right)^d.$$

What concerns the evaluation of $`\text{Fes}(Q)`$, respectively $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, we refer the reader to the next section for the determinant evaluations that we are going to use. In combination with Lemma 5, these will allow us to evaluate Scott-type permanents in numerous special cases, see Section 5.

4. Determinant evaluations

Proposition 6.— Let $`n`$ and $`r`$ be positive integers, $`rn`$, and $`x_1,x_2,\mathrm{},x_n`$, $`y_1,y_2,\mathrm{},y_n`$ be indeterminates. Then, with $`d=\text{gcd}(r,n)`$, we have

$$\begin{array}{cc}& det\left(\begin{array}{ccccccc}x_1& 0& \mathrm{}& 0& y_{nr+1}& 0& \\ 0& x_2& 0& & 0& y_{nr+2}& 0\\ & & \mathrm{}& & & & \mathrm{}& 0\\ 0& & & & & & 0& y_n\\ y_1& 0& \\ 0& y_2& 0\\ & 0& \mathrm{}& 0& & & \mathrm{}& 0\\ & & 0& y_{nr}& 0& & 0& x_n\end{array}\right)\hfill \\ & =\underset{i=1}{\overset{d}{}}\left(\underset{j=1}{\overset{n/d}{}}x_{i+(j1)d}(1)^{n/d}\underset{j=1}{\overset{n/d}{}}y_{i+(j1)d}\right).\hfill \end{array}$$
$`(5)`$
(I.e., in the matrix there are only nonzero entries along two diagonals, one of which is a broken diagonal.)

Proof.—Let $`A=(A_{ij})_{1i,jn}`$ be the matrix of which the determinant is taken in (5). By definition, we have

$$detA=\underset{\sigma S_n}{}\text{sgn}\sigma \underset{i=1}{\overset{n}{}}A_{i\sigma (i)},$$
$`(6)`$
where $`S_n`$ denotes the symmetric group of order $`n`$. Since the matrix $`A`$ is a sparse matrix, only those permutations $`\sigma `$ contribute to the sum which have the property $`\sigma (i)=i`$ or $`\sigma (i)=i+r`$ mod $`n`$ for all $`i`$. Thus, the decomposition into disjoint cycles of such a permutation consists only of cycles of length $`1`$, and of cycles of length $`n/d`$ of the form $`(i,i+r,\mathrm{},i+(n/d1)r)`$ (where, again, all integers have to be taken modulo $`n`$). Using these observations and the fact that the sign of any cycle of length $`n/d`$ is $`(1)^{n/d1}`$ in (6) yields (5) immediately.

Theorem 7.— Let $`a,b,c,d,e`$ be indeterminates. For any positive integer $`n`$ and integers $`i,j`$ let $`n(i,j)`$ denote 1 plus (the representative between $`0`$ and $`n1`$ of) the residue class of $`ij`$ mod $`n`$. Then

$$\begin{array}{c}\text{ }\underset{1i,jn}{det}\left((n(i,j)+c)(n(i,j)a+b)+d(j1)(n(i,j)a+e)\right)\text{ }\hfill \\ \text{ }=det\left(\begin{array}{cccc}\left(\mathrm{c}+1\right)\left(\mathrm{a}+\mathrm{b}\right)& \left(\mathrm{c}+\mathrm{n}\right)\left(\mathrm{na}+\mathrm{b}\right)& & \left(\mathrm{c}+2\right)\left(2\mathrm{a}+\mathrm{b}\right)\\ +\mathrm{d}& +\mathrm{d}\left(\mathrm{na}+\mathrm{e}\right)& \mathrm{}& +\mathrm{d}\left(\mathrm{n}1\right)\left(2\mathrm{a}+\mathrm{e}\right)\\ \\ \left(\mathrm{c}+2\right)\left(2\mathrm{a}+\mathrm{b}\right)& \left(\mathrm{c}+1\right)\left(\mathrm{a}+\mathrm{b}\right)& & \left(\mathrm{c}+3\right)\left(3\mathrm{a}+\mathrm{b}\right)\\ +\mathrm{d}& +\mathrm{d}\left(\mathrm{a}+\mathrm{e}\right)& \mathrm{}& +\mathrm{d}\left(\mathrm{n}1\right)\left(3\mathrm{a}+\mathrm{e}\right)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \\ \left(\mathrm{c}+\mathrm{n}\right)\left(\mathrm{na}+\mathrm{b}\right)& \left(\mathrm{c}+\mathrm{n}1\right)& & \left(\mathrm{c}+1\right)\left(\mathrm{a}+\mathrm{b}\right)\\ +\mathrm{d}& \times \left(\left(\mathrm{n}1\right)\mathrm{a}+\mathrm{b}\right)+\mathrm{d}& \mathrm{}& +\mathrm{d}\left(\mathrm{n}1\right)\left(\mathrm{a}+\mathrm{e}\right)\\ & \left(\left(\mathrm{n}1\right)\mathrm{a}+\mathrm{e}\right)& & \end{array}\right)\text{ }\hfill \\ \text{ }=(n)^{n1}U_n(a,b,c,d,e)\underset{i=3}{\overset{n}{}}(ia+b+ca),(7)\text{ }\hfill \end{array}$$

where $`U_n(a,b,c,d,e)`$ is the polynomial

$$\begin{array}{cc}& \mathrm{U}_\mathrm{n}(\mathrm{a},\mathrm{b},\mathrm{c},\mathrm{d},\mathrm{e})=\frac{\left(\mathrm{n}+1\right)\left(\mathrm{n}+2\right)}{3}\mathrm{a}^2+\frac{\left(\mathrm{n}+1\right)\left(2\mathrm{n}+7\right)}{6}\mathrm{ab}+\frac{\left(\mathrm{n}+1\right)}{2}\mathrm{b}^2\hfill \\ & +\frac{\left(\mathrm{n}+1\right)\left(2\mathrm{n}+7\right)}{6}\mathrm{a}^2\mathrm{c}+\frac{\left(3\mathrm{n}+5\right)}{2}\mathrm{abc}+\mathrm{b}^2\mathrm{c}+\frac{\left(\mathrm{n}+1\right)}{2}\mathrm{a}^2\mathrm{c}^2+\mathrm{abc}^2+\frac{\left(\mathrm{n}+3\right)}{2}\mathrm{ad}\hfill \\ & +\mathrm{bd}+\mathrm{acd}\frac{\left(\mathrm{n}1\right)\left(2\mathrm{n}+5\right)}{6}\mathrm{ae}\frac{\left(\mathrm{n}1\right)}{2}\mathrm{be}\frac{\left(\mathrm{n}1\right)}{2}\mathrm{ace}.\hfill \end{array}$$

In the case that $`n=2`$, the product in (7) has to be read as $`1`$, and in the case that $`n=1`$, the product has to be interpreted as $`1/(2a+b+ca)`$.

Proof.—In the cases $`n=1`$ and $`n=2`$, the claim can be verified directly. For $`n3`$, we use the “identification of factors” method as explained in \[10, Sec. 2.4\] or \[9, Sec. 2\].

We proceed in several steps. An outline is as follows. In the first step we show that $`_{i=3}^n(ia+b+ca)`$ is a factor of the determinant as a polynomial in $`a,b,c,d,e`$. In the second step we prove that $`U_n(a,b,c,d,e)`$ is a factor of the determinant. Then, in the third step, we determine the maximal degree of the determinant as a polynomial in $`a`$, and also in $`b`$, $`c`$, $`d`$, and in $`e`$. It turns out that the maximal degree is $`n`$ as a polynomial in $`a`$, the same being true as a polynomial in $`b`$ and as a polynomial in $`c`$, while it is 1 as a polynomial in $`d`$, the same being true as a polynomial in $`e`$. On the other hand, the degree in $`a`$, and also in $`b`$ and in $`c`$, of the product on the right-hand side of (7), which by the first two steps divides the determinant, is exactly $`n`$. It is exactly 1 in $`d`$ and also in $`e`$. Therefore we are forced to conclude that the determinant equals

$$C(n)U_n(a,b,c,d,e)\underset{i=3}{\overset{n}{}}(ia+b+ca),$$
$`(8)`$
where $`C(n)`$ is a constant independent of $`a,b,c,d,e`$. Finally, in the fourth step, we determine the constant $`C(n)`$, which turns out to equal $`(n)^{n1}`$. Clearly, this would finish the proof of theorem.

Step 1. For $`i=3,\mathrm{},n`$ the term $`(ia+b+ca)`$ is a factor of the determinant. We claim that, if $`b=iaca`$, we have

$$\begin{array}{cc}& \left(\text{row }(ni)\right)+3\left(\text{row }(ni+1)\right)\hfill \\ & 3\left(\text{row }(ni+2)\right)+\left(\text{row }(ni+3)\right)=0\hfill \end{array}$$

as long as $`n>i`$. (Here, (row $`i`$) denotes the $`i`$-th row of the matrix underlying the determinant in (7).) In the case that $`n=i`$, we claim that we have

$$3\left(\text{row }1\right)+3\left(\text{row }2\right)\left(\text{row }3\right)+\left(\text{row }n\right)=0$$

as long as $`n>3`$, and that we have $`\left(\text{row }1\right)+\left(\text{row }2\right)=0`$ if $`n=3`$. All these claims are easily verified by an obvious case-by-case analysis.

Step 2. The polynomial $`U_n(a,b,c,d,e)`$ is a factor of the determinant. We claim that if $`d`$ is chosen so that $`U_n(a,b,c,d,e)`$ vanishes, we have

$$\underset{j=1}{\overset{n}{}}((j+1)a+b+ca)\left(\text{column }j\right)=0.$$

Again, it is a routine task to verify this identity.

Step 3. The determinant is a polynomial in $`a`$ (in $`b`$, respectively in $`c`$) of maximal degree $`n`$, and a polynomial in $`d`$ (respectively in $`e`$) of maximal degree $`1`$. The first claim follows from the fact that each term in the defining expansion of the determinant has degree $`n`$ in $`a`$ (as well as in $`b`$, respectively in $`c`$). To establish the second claim, we simply subtract the first row of the determinant from all other rows, with the effect that only the entries in the first row contain $`d`$ and $`e`$ after these transformations. Since the right-hand side of (7), which by Steps 1 and 2 divides the determinant as a polynomial in $`a,b,c,d,e`$, also has degree $`n`$ in $`a`$, in $`b`$, and in $`c`$, and degree $`1`$ in $`d`$, and in $`e`$, the determinant and the right-hand side of (7) differ only by a multiplicative constant.

Step 4. The evaluation of the multiplicative constant. By the preceding steps we know that the determinant equals (8). In particular, if we set $`a=c=d=e=0`$ and $`b=1`$, we have

$$\underset{1i,jn}{det}\left(n(i,j)\right)=C(n)(n+1)/2.$$
$`(9)`$
The matrix on the left-hand side of (9) is a circulant matrix with entries $`1,2,\mathrm{},n`$. Hence, its determinant equals

$$\underset{\omega \text{ : zero of }x^n1}{}(1+2\omega +3\omega ^2+\mathrm{}+n\omega ^{n1}).$$

The sum is easily evaluated by observing that it is the derivative of a geometric series. It turns out to be equal to $`n/(1\omega )`$. The resulting product simplifies by the observation

$$\underset{\omega \text{ : zero of }x^n1,\omega 1}{}(1\omega )=(1+x+\mathrm{}+x^{n1})|_{x=1}=n.$$
$`(10)`$
Thus, the determinant in (9) equals $`(n)^{n1}(n+1)/2`$. Therefore $`C(n)`$ is equal to $`(n)^{n1}`$.

This finishes the proof of (7) and thus of the theorem.

Theorem 8.— Let $`a`$ be an indeterminate. For any positive integer $`n`$ and integers $`i,j`$ let $`s(i,j)`$ denote (the representative between $`0`$ and $`n1`$ of) the residue class of $`ij+1`$ mod $`n`$. Then

$$\begin{array}{c}\text{ }\underset{1i,jn1}{det}\left(\left\{\begin{array}{cc}\left(\mathrm{n}\mathrm{m}1\right)+\mathrm{j}\left(1+\mathrm{a}\mathrm{n}\right)& \text{if }\mathrm{i}=\mathrm{j}2\left(\text{mod }\mathrm{n}\right)\\ \\ \left(\mathrm{n}1\right)\left(\mathrm{m}1\right)+\mathrm{j}\left(1\mathrm{a}\right)& \text{if }\mathrm{i}=\mathrm{j}3\left(\text{mod }\mathrm{n}\right)\\ \\ \left(\mathrm{n}\mathrm{m}32\mathrm{s}(\mathrm{i},\mathrm{j})\right)+\mathrm{j}& \text{otherwise}\end{array}\right\}\right)\text{ }\hfill \\ \text{ }=(1)^{n1}\frac{1}{n}\underset{i=2}{\overset{n}{}}(nmia).(11)\text{ }\hfill \end{array}$$

Proof.—The matrix underlying this determinant is a matrix whose elements have a uniform definition, except for two (broken) diagonals, the one with $`i=j2(\text{mod }n)`$, and the one with $`i=j3(\text{mod }n)`$.

To begin with, we reorder the rows so that the next-to-last row becomes the first row, the last row becomes the second row, and then follow the remaining rows in their original order.

Now we add all the rows to the (new) last row. In the resulting matrix, we change the sign of the last row and, subsequently, move it up so that it becomes the third row. As a result of these manipulations, we obtain

$$(1)^{n1}\underset{1i,jn1}{det}\left(\left\{\begin{array}{cc}(n1)(m1)+j(1a)& \text{if }i=j\\ (nm1)+j(an+1)& \text{if }i=j+1\\ mn+12i+3j& \text{if }i<j\\ m+n+12i+3j& \text{if }i>j+1\end{array}\right\}\right).$$
$`(12)`$
Next we apply further row operations. We subtract the second row from the first, the third from the second, …, the $`(n1)`$-st row from the $`(n2)`$-nd row, in that order. Subsequently, we repeat the same kind of operations, but stop before the last row, i.e., we subtract the second row from the first, the third from the second, …, the $`(n2)`$-nd row from the $`(n3)`$-rd row, in that order. As a result, the above determinant is converted into the determinant of the following matrix

$$\left(\left\{\begin{array}{cc}nmja& \text{if }i=j2\\ 2nm+2nj(3an)& \text{if }i=j1\text{ and }i<n2\\ nm2nj(3a2n)& \text{if }i=j\text{ and }i<n2\\ j(an)& \text{if }i=j+1\text{ and }i<n2\\ 2& \text{if }i=n2\text{ and }jn4\\ (n3)(an)+2& \text{if }i=n2\text{ and }j=n3\\ nm2n+2(n2)(2an)& \text{if }i=n2\text{ and }j=n2\\ nnm+1+(n1)(a1)& \text{if }i=n2\text{ and }j=n1\\ mn+3+3j& \text{if }i=n1\text{ and }jn3\\ nm1+(n2)(an+1)& \text{if }i=n1\text{ and }j=n2\\ (n1)(ma)& \text{if }i=n1\text{ and }j=n1\\ 0& \text{otherwise}\end{array}\right\}\right).$$
$`(13)`$
This is a matrix with four “special” diagonals (the diagonals with $`i=j2`$, $`j1`$, $`j`$, $`j+1`$, respectively) and two “special” rows (the last two rows). All other entries are zero.

Now we factor $`(nm3a)`$ out of the first row. Subsequently, we substract $`(an)`$ times the (new) first row from the second. Now one is able to factor $`(nm4a)`$ out of the (new) second row. Next, we substract $`2(an)`$ times the (new) second row from the third row. Etc. We stop this procedure in the $`(n3)`$-rd row. Thus, the determinant of the matrix in (13) is equal to $`_{i=3}^{n1}(nmia)`$ times the determinant of the matrix

$$\left(\begin{array}{ccccc}1& 2& 1& 0& \mathrm{}\\ 0& 1& 2& 1& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ & & & 0& 1& 2& 1\\ a_{n2,1}& a_{n2,2}& \multicolumn{4}{c}{}& a_{n2,n1}\\ a_{n1,1}& a_{n1,2}& \multicolumn{4}{c}{}& a_{n1,n1}\end{array}\right),$$
$`(14)`$
where the entries in the last two rows are still the same as in (13).

We now perform the final set of transformations. We add column 1 through column $`n2`$ to column $`n1`$, and then we add

$$\underset{j=1}{\overset{n3}{}}(nj1)\left(\text{column }j\right)$$

to column $`n2`$. The effect is that a block matrix is obtained of the form

$$\left(\begin{array}{cc}U& 0\\ & M\end{array}\right),$$

where $`U`$ is an $`(n3)\times (n3)`$ upper-triangular matrix with 1s on the diagonal, and where $`M`$ is the $`2\times 2`$ matrix

$$\left(\begin{array}{cc}nmn2a+2& n2\\ \frac{1}{2}(n2)(n+m+2anm3)& \frac{5}{2}n\frac{n^2}{2}+ma3\end{array}\right).$$

Clearly, the determinant of $`U`$ is 1, while the determinant of $`M`$ is $`(ma)(nm2a)`$. Putting everything together, we have completed the proof of (11).

Corollary 9.— Let $`a`$ be an indeterminate. For any positive integer $`n`$ and integers $`i,j`$ let $`s(i,j)`$ denote (the representative between $`0`$ and $`n1`$ of) the residue class of $`ij+1`$ mod $`n`$. Then

$$\begin{array}{c}\text{ }\underset{1i,jn1}{det}\left(\left\{\begin{array}{cc}2s(i,j)a+j1& \text{if }ij2(\text{mod }n)\\ (n1)(n+aj1)& \text{if }i=j2(\text{mod }n)\end{array}\right\}\right)\text{ }\hfill \\ \text{ }=n^{n2}\underset{i=0}{\overset{n2}{}}(i+a).(15)\text{ }\hfill \end{array}$$

Proof.—In the determinant, we move the last row on top, replace the (now) last row by the sum of all the rows, factor $`(1)`$ out of the resulting row, and finally move it up so that it becomes the second row, retaining the order of all the other rows. These operations did not change the value of the determinant. However, the resulting determinant is exactly the determinant in (12) with $`a=n`$ and $`m=a+n`$. As we have shown in the proof of Theorem 8, the latter determinant differs from the determinant in (11) just by a sign of $`(1)^{n1}`$. Thus, we obtain the right-hand side of (15).

5. Closed form evaluations for Scott-type permanents

Theorem 10.— Let $`n`$, $`m`$ and $`r`$ be positive integers and $`d=\text{gcd}(n,r)`$. Then

$$\begin{array}{cc}& \text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}y^\mathrm{}n+\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}y^{\mathrm{}n+r})\hfill \\ & =\frac{\begin{array}{cc}& d^n{\displaystyle \underset{i=1}{\overset{d}{}}}(\left({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}a_{\mathrm{}}\right)^{n/d}({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}(in\mathrm{}1)a_{\mathrm{}}/{\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}da_{\mathrm{}})_{n/d}\hfill \\ & ({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}b_{\mathrm{}})^{n/d}({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}(irn\mathrm{}1)b_{\mathrm{}}/{\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}db_{\mathrm{}})_{n/d})\hfill \end{array}}{\left(\left({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}a_{\mathrm{}}\right)^{n/d}\left({\displaystyle \underset{\mathrm{}=0}{\overset{m}{}}}b_{\mathrm{}}\right)^{n/d}\right)^d},\hfill \end{array}$$

where $`(\alpha )_k`$ is the standard notation for shifted factorials, $`(\alpha )_k:=\alpha (\alpha +1)\mathrm{}(\alpha +k1)`$, $`k1`$, and $`(\alpha )_0:=1`$.

Proof.—Let us first consider the case that $`nr`$. According to Theorem 3, we have to compute the quotient $`\text{Fes}(Q)/\text{Res}(x^n1,Q)`$, where $`Q=_{\mathrm{}=0}^ma_{\mathrm{}}y^\mathrm{}n+_{\mathrm{}=0}^mb_{\mathrm{}}y^{\mathrm{}n+r}`$. In order to compute $`\text{Fes}(Q)`$, in (3) we replace $`a_\mathrm{}n`$ by $`a_{\mathrm{}}`$ and $`a_{\mathrm{}n+r}`$ by $`b_{\mathrm{}}`$, $`\mathrm{}=0,1,\mathrm{},m`$, and set all other $`a_{\mathrm{}}`$’s equal to zero. In the resulting determinant we move the last row to the top, thus creating a sign of $`(1)^{n1}`$, and finally apply Proposition 6 with $`x_j=_{\mathrm{}=0}^m(n\mathrm{}j+1)a_{\mathrm{}}`$ and $`y_j=_{\mathrm{}=0}^m(n\mathrm{}+rj+1)b_{\mathrm{}}`$, $`j=1,2,\mathrm{},n`$. For the evaluation of $`\text{Res}(x^n1,Q)`$ we note that

$$\begin{array}{cc}& \text{Res}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}y^\mathrm{}n+\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}y^{\mathrm{}n+r})\hfill \\ & =\underset{\omega \text{ : zero of }x^n1}{}\left(\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}+\omega ^r\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}\right)=\text{Res}(x^n1,y^r\underset{\mathrm{}=0}{\overset{m}{}}b_{\mathrm{}}+\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}).\hfill \end{array}$$

Now we can apply Lemma 5.

If on the other hand $`nr`$, then $`d=n`$. It can be verified directly that the claimed formula remains valid in that case, too.

Remark.— Given a polynomial $`Q(y)`$, there is no unique way to write it in the form $`_{\mathrm{}=0}^ma_{\mathrm{}}y^\mathrm{}n+_{\mathrm{}=0}^mb_{\mathrm{}}y^{\mathrm{}n+r}`$. For example, we may write $`Q(y)=y^n+a+b`$ as $`Q(y)=(y^n+ay^0)+b`$, or as $`Q(y)=(y^n+(a+b))`$ (i.e., either with $`a_1=1`$, $`a_0=a`$, $`r=0`$, $`b_0=b`$, or with $`a_1=1`$, $`a_0=a+b`$, $`b_{\mathrm{}}=0`$ for all $`\mathrm{}`$). Regardless which choice we make, Theorem 10 yields

$$\text{PER}(x^n1,y^n+a+b)=(1)^{n+1}\frac{\underset{i=1}{\overset{n}{}}(i(ni)(a+b))}{(a+b+1)^n}.$$

Corollary 11.— Let $`n`$ and $`m`$ be positive integers. Then

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}y^\mathrm{}n)=\left(n\underset{\mathrm{}=0}{\overset{m}{}}\mathrm{}a_{\mathrm{}}/\underset{\mathrm{}=0}{\overset{m}{}}a_{\mathrm{}}\right)_n.$$

Corollary 12.— We have

$$\text{PER}(x^n1,y^{mn}+\mathrm{}+y^{2n}+y^n+1)=(mn/2)_n.$$

Corollary 13.— If $`m`$ is even then

$$\text{PER}(x^n+1,y^{mn}+\mathrm{}+y^{2n}+y^n+1)=(mn/2)_n.$$

Proof.—We use the case $`a_{\mathrm{}}=(1)^{\mathrm{}}`$ in Corollary 11, and the fact that

$$\text{per}\left(\frac{1}{x_i\sqrt[n]{1}y_j}\right)_{1in,\mathrm{\hspace{0.17em}1}jmn}=\text{per}\left(\frac{1}{x_iy_j/\sqrt[n]{1}}\right)_{1in,\mathrm{\hspace{0.17em}1}jmn}.$$

Corollary 14.— We have

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}\mathrm{}y^\mathrm{}n)=(n(2m+1)/3)_n.$$

Corollary 15.— We have

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{m}{}}\mathrm{}y^{\mathrm{}^2n})=(nm(m+1)/2)_n.$$

Corollary 16.— We have

$$\text{PER}(x^n1,y^{mn}+ay^{rn}+b)=((m+ra)n/(a+b+1))_n.$$

Corollary 17.— We have

$$\text{PER}(x^n1,y^{mn}+1)=(mn/2)_n.$$

For $`m=1`$ one recovers Scott’s identity stated at the beginning of the introduction.

Corollary 18.— If $`m+ra=a+b+10`$, then

$$\text{PER}(x^n1,y^{mn}+ay^{rn}+b)=(1)^{n+1}n!.$$

Corollary 19.— If $`a2`$ then

$$\text{PER}(x^n1,y^{2n}+ay^n+1)=(1)^{n+1}n!.$$

Corollary 20.— If $`m+ra=0`$ and $`a+b+10`$, then

$$\text{PER}(x^n1,y^{mn}+ay^{rn}+b)=0.$$

Corollary 21.— If $`b1`$ then

$$\text{PER}(x^n1,y^{2n}2y^n+b)=0.$$

Corollary 22.— Let $`n`$ and $`m`$ be positive integers and $`d=\text{gcd}(n,m)`$. Then

$$\text{PER}(x^n1,y^m+b)=\frac{d^n{\displaystyle \underset{i=1}{\overset{d}{}}}\left(\left(\frac{im1}{d}\right)_{n/d}(b)^{n/d}\left(\frac{i1}{d}\right)_{n/d}\right)}{\left(1(b)^{n/d}\right)^d}.$$

Proof.—In Theorem 10, set $`m=0`$, $`a_0=b`$, $`b_0=1`$, and replace $`r`$ by $`m`$, in this order.

Corollary 23.— If $`\text{gcd}(m,n)=1`$, then

$$\text{PER}(x^n1,y^m+b)=(1)^{n+1}\frac{m(m1)\mathrm{}(mn+1)}{1(b)^n}.$$

Corollary 24.— If $`\text{gcd}(m,n)=1`$, then

$$\begin{array}{cc}& \text{PER}(x^{s(n1)}+\mathrm{}+x^{2s}+x^s+1,y^{s(m1)}+\mathrm{}+y^{2s}+y^s+1)\hfill \\ & =\frac{{\displaystyle \underset{i=0}{\overset{s1}{}}}\left({\displaystyle \underset{\mathrm{}=0}{\overset{n1}{}}}(i+\mathrm{}s){\displaystyle \underset{\mathrm{}=0}{\overset{n1}{}}}(i+\mathrm{}sms)\right)}{(mns)^s}.\hfill \end{array}$$

Proof.—Consider the Scott-type permanent $`\text{PER}(x^{sn}1,y^{sm}q^{sm})`$ (expressed in terms of its definition). When we multiply it by $`(1q)^s`$ and then perform the limit $`q1`$, then the permanent reduces to $`\text{PER}(x^{s(n1)}+\mathrm{}+x^s+1,y^{s(m1)}+\mathrm{}+y^s+1)`$, as is straightforward to see. On the other hand, the permanent that we started with is the permanent in Corollary 22 with $`n`$ replaced by $`sn`$, $`m`$ replaced by $`sm`$, and $`b=q^{sm}`$. Indeed, if we multiply the right-hand side from Corollary 22 (with these choices for the parameters) by $`(1q)^{sm}`$, and then perform the limit $`q1`$, we obtain exactly the claimed result.

Corollary 25.— If $`\text{gcd}(m,n)=1`$, then

$$\begin{array}{cc}& \text{PER}(x^{n1}+\mathrm{}+x+1,y^{m1}+\mathrm{}+y+1)\hfill \\ & =(1)^{n+1}\frac{(m1)\mathrm{}(mn+1)}{n}.\hfill \end{array}$$

Corollary 26.— If $`\text{gcd}(m,n)=1`$ and $`n`$ is odd, then

$$\text{PER}(x^n1,y^m+1)=\frac{m(m1)\mathrm{}(mn+1)}{2}.$$

Corollary 27.— If $`n`$ is odd, then

$$\text{PER}(x^n1,y^{n+1}+1)=\frac{(n+1)!}{2}.$$

Corollary 28.— Let $`n`$ and $`r`$ be positive integers (not necessarily $`n>r`$) and $`d=\text{gcd}(n,r)`$. Then

$$\begin{array}{cc}& \text{PER}(x^n1,y^n+ay^r+b)\hfill \\ & =\frac{d^n{\displaystyle \underset{i=1}{\overset{d}{}}}\left((b+1)^{\frac{n}{d}}\left({\displaystyle \frac{ibb+in1}{d(b+1)}}\right)_{\frac{n}{d}}(a)^{\frac{n}{d}}\left({\displaystyle \frac{ir1}{d}}\right)_{\frac{n}{d}}\right)}{\left((b+1)^{\frac{n}{d}}(a)^{\frac{n}{d}}\right)^d}.\hfill \end{array}$$

Proof.—In Theorem 10, set $`m=1`$, $`a_0=b`$, $`a_1=1`$, $`b_0=a`$, and $`b_1=0`$.

Corollary 29.— Let $`n`$ and $`r`$ be positive integers (not necessarily $`n>r`$) and $`\text{gcd}(n,r)=1`$. Then

$$\text{PER}(x^n1,y^n+ay^r+b)=(1)^{n+1}\frac{\underset{i=1}{\overset{n}{}}(i(ni)b))a^n(r)_n}{(b+1)^n(a)^n}.$$

If $`1rn1`$ then $`(r)_n=0`$. Thus, one recovers the results in . On the other hand, if we set $`r=n+1`$, we obtain, for example, the following result.

Corollary 30.— We have

$$\text{PER}(x^n1,y^{n+1}+y^n1)=n^n(1)^n(n+1)!.$$

Corollary 31.— We have

$$\text{PER}(x^n1,y^n+ny1)=1.$$

Theorem 32.— Let $`n`$ and $`m`$ be positive integers and $`a`$ be an arbitrary number. Then

$$\begin{array}{cc}& \text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+a)y^{\mathrm{}})\hfill \\ & =(1)^{n1}\frac{n(m1)V_n(a,m)}{6\left(mn+2a1\right)}\left(a+(m1)n+1\right)_{n2},\hfill \end{array}$$

where $`V_n(a,m)`$ is the polynomial

$$V_n(a,m)=16a+6a^2+n2an5mn+10amnmn^2+4m^2n^2.$$

Proof.—Let $`Q=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^{\mathrm{}}`$. Using Theorem 3 again, we have to compute $`\text{Fes}(Q)/\text{Res}(x^n1,Q)`$. In order to compute $`\text{Fes}(Q)`$, in (3) set $`a_{\mathrm{}}=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn1`$. In the resulting determinant we move the last row to the top, thus creating a sign of $`(1)^{n1}`$, and finally apply Theorem 7 with $`a=1`$, $`b=(m1)n1`$, $`c=a1`$, $`d=n^2(m1)(2m1)/6an(m1)/2`$ and $`e=a+n(m1)/21`$.

For the computation of $`\text{Res}(x^n1,Q)`$ we note that

$$\begin{array}{cc}& \text{Res}(x^n1,Q)=\underset{\omega \text{ : zero of }x^n1}{}\left(\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+a)\omega ^{\mathrm{}}\right)\hfill \\ & =\left(\left(\genfrac{}{}{0pt}{}{mn}{2}\right)+mna\right)\underset{\omega \text{ : zero of }x^n1,\omega 1}{}\left(\underset{\mathrm{}=1}{\overset{mn}{}}\mathrm{}\omega ^\mathrm{}1\right).\hfill \end{array}$$

The sum in the last line is the derivative of a geometric series, and is therefore easily evaluated. The result of the summation turns out to be $`mn/(1\omega )`$. The computation is completed by the observation (10), and some simplification.

Corollary 33.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}\mathrm{}y^{\mathrm{}})=(1)^{n1}\frac{(4mnn1)(mn2)!}{6(mnn1)!}.$$

Corollary 34.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+1)y^{\mathrm{}})=(1)^{n1}\frac{(4mnn+1)(mnn)(mn1)!}{6(mnn+1)!}.$$

Corollary 35.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(mn\mathrm{})y^{\mathrm{}})=\frac{(m1)(n+1)!}{6}.$$

Corollary 36.— Let $`n`$ and $`m`$ be positive integers, $`n2`$. Then

$$\text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{mn1}{}}(mn\mathrm{}1)y^{\mathrm{}})=\frac{(m1)n!}{6}.$$

Remark.— It is also possible to move forward and derive formulas for $`\text{PER}(x^n1,_{\mathrm{}=0}^{mn/s1}(\mathrm{}+a)y^\mathrm{}s)`$, where $`s`$ is some positive integer. This would require to find analogues of Theorem 7 in which $`n(i,j)`$ is replaced by $`Dn(i,j)`$, where $`D`$ is the inverse of $`s/\text{gcd}(n,s)`$ modulo $`n/\text{gcd}(n,s)`$. As calculations aided by the computer indicate, the resulting determinant evaluations have forms very similar to (7). That is, the result shows a product of linear factors in $`a`$ and $`b`$ as the one on the right-hand side of (7), and one irreducible polynomial of higher degree (such as $`U_n(a,b,c,d,e)`$ in (7)). However, as $`D`$ increases, the degree of the irreducible polynomial also increases, whereas the amount of linear factors in $`a`$ and $`b`$ decreases, so that the results become more and more unwieldy. We therefore content ourselves with stating the result when $`s`$ divides $`n`$.

Theorem 37.— Let $`n`$, $`m`$ and $`s`$ be positive integers so that $`sn`$, and let $`a`$ be an arbitrary number. Then

$$\begin{array}{cc}& \text{PER}(x^n1,\underset{\mathrm{}=0}{\overset{\frac{mn}{s}1}{}}(\mathrm{}+a)y^\mathrm{}s)=(1)^{n1}\frac{s^{n2s}}{6^s(mn+2ass)^s}\hfill \\ & \times \underset{k=0}{\overset{s1}{}}\left((a+\frac{1}{s}(nmnk)+1)_{n/s2}V_{n,s}(a,m,k)\right),\hfill \end{array}$$

where $`V_{n,s}(a,m,k)`$ is the polynomial

$$\begin{array}{cc}\hfill \mathrm{V}_{\mathrm{n},\mathrm{s}}& (\mathrm{a},\mathrm{m},\mathrm{k})=6\mathrm{k}^2\mathrm{mn}+6\mathrm{k}\mathrm{m}\mathrm{n}^210\mathrm{k}\mathrm{m}^2\mathrm{n}^2+\mathrm{mn}^35\mathrm{m}^2\mathrm{n}^3+4\mathrm{m}^3\mathrm{n}^36\mathrm{k}^2\mathrm{s}\hfill \\ & +12\mathrm{a}\mathrm{k}^2\mathrm{s}6\mathrm{k}\mathrm{n}\mathrm{s}+12\mathrm{a}\mathrm{k}\mathrm{n}\mathrm{s}+12\mathrm{k}\mathrm{m}\mathrm{n}\mathrm{s}24\mathrm{a}\mathrm{k}\mathrm{m}\mathrm{n}\mathrm{s}\mathrm{n}^2\mathrm{s}+2\mathrm{a}\mathrm{n}^2\mathrm{s}+6\mathrm{m}\mathrm{n}^2\mathrm{s}\hfill \\ & 12\mathrm{a}\mathrm{m}\mathrm{n}^2\mathrm{s}5\mathrm{m}^2\mathrm{n}^2\mathrm{s}+10\mathrm{a}\mathrm{m}^2\mathrm{n}^2\mathrm{s}2\mathrm{k}\mathrm{s}^2+12\mathrm{a}\mathrm{k}\mathrm{s}^212\mathrm{a}^2\mathrm{ks}^2\mathrm{ns}^2\hfill \\ & +6\mathrm{a}\mathrm{n}\mathrm{s}^26\mathrm{a}^2\mathrm{ns}^2+\mathrm{mns}^26\mathrm{a}\mathrm{m}\mathrm{n}\mathrm{s}^2+6\mathrm{a}^2\mathrm{mns}^2.\hfill \end{array}$$

Proof.—Let $`Q=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^\mathrm{}s`$. Again, according to Theorem 3, we have to compute $`\text{Fes}(Q)/\text{Res}(x^n1,Q)`$. In order to compute $`\text{Fes}(Q)`$, in (3) set $`a_\mathrm{}s=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn/s1`$, and all other $`a_i`$’s to zero. In the resulting determinant we move the last row to the top, thus creating a sign of $`(1)^{n1}`$. Since only every $`s`$-th $`a_i`$ is nonzero, we are dealing with a determinant of a matrix in which a lot of entries are zero. If we permute rows and columns so that first come the rows and columns whose indices are congruent $`1`$ mod $`s`$, then come the rows and columns whose indices are congruent $`2`$ mod $`s`$, etc., then the matrix of which we want to compute the determinant assumes a block form, with $`(n/s)\times (n/s)`$ blocks on the diagonal, and zeroes otherwise. Therefore the determinant equals the product of the determinants of the $`s`$ matrices of dimension $`(n/s)\times (n/s)`$ on the diagonal. Each of these determinants can be evaluated by means of Theorem 7. We leave it to the reader to fill in the details.

For the computation of $`\text{Res}(x^n1,Q)`$ we proceed as in the proof of Theorem 32.

Theorem 38.— Let $`n`$ and $`m`$ be positive integers, $`n2`$, and $`a`$ be an arbitrary number. Then

$$\begin{array}{cc}& \text{PER}(x^{n1}+\mathrm{}+x+1,\underset{\mathrm{}=0}{\overset{mn1}{}}(\mathrm{}+a)y^{\mathrm{}})\hfill \\ & =(1)^{n1}\left(a+(m1)n+1\right)_{n1}.\hfill \end{array}$$

Proof.—Let $`P(x)=x^{n1}+\mathrm{}+x+1`$ and $`Q(y)=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^{\mathrm{}}`$. This time we apply Theorem 4. According to that theorem, we have to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)/\text{Res}(P,Q)`$. In order to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, in (4) we set $`a_{\mathrm{}}=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn1`$. The resulting determinant is exactly $`m^{n1}`$ times the determinant in Corollary 9 with $`a`$ replaced by $`((m1)n+a+1)`$. For the computation of the resultant of $`P`$ and $`Q`$ we proceed as in the proof of Theorem 7. Simplification of the result yields the claimed expression.

Theorem 39.— Let $`n`$ and $`m`$ be positive integers, $`n2`$, and $`a`$ be an arbitrary number. Then

$$\begin{array}{cc}& \text{PER}(x^{n1}+\mathrm{}+x+1,\underset{\mathrm{}=0}{\overset{mn2}{}}(\mathrm{}+a)y^{\mathrm{}})\hfill \\ & =(1)^{n1}\frac{(nmn)_{n1}(nm+a1)^{n1}}{(mn+a1)^n(a1)^n}.\hfill \end{array}$$

Proof.—Let $`P(x)=x^{n1}+\mathrm{}+x+1`$ and $`Q(y)=_{\mathrm{}=0}^{mn1}(\mathrm{}+a)y^{\mathrm{}}`$. Using Theorem 4 again, we have to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)/\text{Res}(P,Q)`$. In order to compute $`\stackrel{𝑒}{\mathrm{Fes}}(Q)`$, in (4) we set $`a_{\mathrm{}}=\mathrm{}+a`$, $`\mathrm{}=0,1,\mathrm{},mn2`$. The resulting determinant is exactly $`m^{n1}`$ times the determinant in Theorem 8 with $`m`$ replaced by $`mn+A1`$, $`a`$ replaced by $`(mn+A1)/m`$, and $`A`$ replaced by $`a`$, in that order.

For the computation of $`\text{Res}(x^n+\mathrm{}+x+1,Q)`$ we proceed similarly as in the proof of Theorem 32. Using an observation from that proof, we note that

$$\begin{array}{cc}& \text{Res}(x^n+\mathrm{}+x+1,Q)=\underset{\omega \text{ : zero of }x^n1,x1}{}\left(\underset{\mathrm{}=0}{\overset{mn2}{}}(\mathrm{}+a)\omega ^{\mathrm{}}\right)\hfill \\ & =\underset{\omega \text{ : zero of }x^n1,x1}{}\left(\frac{mn}{1\omega }(mn+a1)\omega ^{mn1}\right)\hfill \\ & =(1)^{n1}\underset{\omega \text{ : zero of }x^n1,x1}{}\frac{mn+a1+\omega (1a)}{\omega (1\omega )}\hfill \\ & =\frac{(mn+a1)^n(a1)^n}{n^2m}.\hfill \end{array}$$

The result follows now upon some simplification.

6. Sums of involutions

As in \[6, Sec. 4\], we may obtain interesting summation theorems by combining special evaluations of Scott-type determinants (see Section 5) with Theorem 2. For example, if we combine Corollary 19 with Theorem 2, we obtain the following result.

Proposition 40.—Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then

$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{n+1}{2x_k}=(1)^{n+1}n!,$$

where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.

Similarly, if we combine Corollary 21 with Theorem 2, we obtain the result below.

Proposition 41.—Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then

$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{n1}{2x_k}=0,$$

where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.

If we combine Corollary 31 with Theorem 2, then we obtain the following result.

Proposition 42.—Let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then

$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{2n+(n+1)x_k}{2x_k^2}=1,$$

where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.

Finally, if we combine Corollary 27 with Theorem 2, we obtain the result below.

Proposition 43.—Let $`n`$ be odd, and let $`x_1,x_2,\mathrm{},x_n`$ be the zeroes of $`x^n1`$. Then

$$\underset{\sigma (n)}{}\underset{(ij)\sigma }{}\frac{1}{(x_ix_j)^2}\underset{(k)\sigma }{}\frac{1n+(3+n)x_k}{2(1+x_k)x_k}=\frac{(n+1)!}{2},$$

where the first product is over all transpositions $`(ij)`$ of $`\sigma `$ in its disjoint cycle decomposition, and where the second product is over all fixed points $`k`$ of $`\sigma `$.

References

 G. E. Andrews, I. P. Goulden, D. M. Jackson.—Generalizations of Cauchy’s summation theorem for Schur functions, Trans. Amer. Math. Soc. 310 (1988), pp. 805–820.

 C. W. Borchardt.—Bestimmung der symmetrischen Verbindungen vermittelst ihrer erzeugenden Funktion, Crelle J. 53 (1855), pp. 193–198.

 D. Callan.—On evaluating permanents and a matrix of cotangents, Linear and Multilinear Algebra 38 (1995), pp. 193–205.

 A. L. Cauchy.—Mémoire sur les fonctions alternées et les sommes alternées, Exercices d’analyse et de phys. math. 2 (1841), pp. 151–159.

 G.-N. Han.—Interpolation entre Cauchy et Borchardt, in preparation, 1999.

 G.-N. Han.—Généralisation de l’identité de Scott sur les permanents, Linear Algebra Appl. (to appear).

 J.-P. Jouanolou.—Polynômes cyclotomiques : Théorie élémentaire et applications.—Prépublication, Université de Strasbourg I, 96 pages, 1990.

 R. Kittappa.—Proof of a conjecture of 1881 on permanents, Linear and Multilinear Algebra 10 (1981), pp. 75–82.

 C. Krattenthaler.—An alternative evaluation of the Andrews–Burge determinant, in: Mathematical Essays in Honor of Gian-Carlo Rota, B. E. Sagan, R. P. Stanley, eds., Progress in Math., vol. 161, Birkhäuser, Boston, 1999, pp. 263–270.

 C. Krattenthaler.—Advanced determinant calculus, Séminaire Lotharingien Combin. 42, “The Andrews Festschrift” (1999), paper B42q, 67 pp.

 A. Lascoux.—Square-ice Enumeration, Séminaire Lotharingien Combin. 42, “The Andrews Festschrift” (1999), paper B42p, 15 pp.

 I. G. Macdonald.—Symmetric functions and Hall polynomials, second edition.—Clarendon Press, Oxford, 1995.

 H. Minc.—Permanents.—Encyclopedia of mathematics and its applications, vol. 6, Addison-Wesley, Mass., 1978.

 H. Minc.—On a conjecture of R. F. Scott, Linear Algebra Appl. 28 (1979), pp. 141–153.

 T. Muir.—The theory of determinants in the historical order of development, 4 vols.—Macmillan, London, 19061923.

 R. F. Scott.—Mathematical notes, Messenger of Math. 10 (1881), pp. 142–149.

 D. Svrtan.—Proof of Scott’s conjecture, Proc. Amer. Math. Soc. 87 (1983), pp. 203–207.

I.R.M.A. and C.N.R.S.  Université Louis Pasteur  7, rue René-Descartes  F-67084 Strasbourg, France  guoniu@math.u-strasbg.fr     Institut für Mathematik  Universität Wien  Strudlhofgasse 4  A-1090 Vienna, Austria  kratt@ap.univie.ac.at


text:
# Kinetics of Ordering in Fluctuation-Driven First-Order Transitions: Simulations and Dynamical Renormalization

## I Introduction

Kinetics of growth can be broadly classified into two categories; nucleation and spinodal decomposition (or continuous ordering). The former applies to situations where the initial state is metastable and the latter to where the initial state is unstable. The distinction becomes unclear near the limit of metastability and in systems where the concept of metastability itself is ill-defined. A class of systems where the definition of metastability becomes ambiguous is one where there is a fluctuation-driven first-order phase transition. A notable example of a system undergoing such a transition is a symmetric di-block copolymer. These systems undergo microphase separation and there is a temperature-driven transition from a uniform phase to a lamellar phase. Another system which is governed by the same dynamical equation as block copolymers, and so exhibits a similar transition, is a Rayleigh-Benard Cell .

A theoretical model describing these transitions was proposed by Brazovskii in 1975. Di-block copolymers were shown to belong to the Brazovskii class by Leibler and the theoretical predictions regarding the equilibrium nature of the ordering transition have been experimentally verified. It was shown by Brazovskii, within a self-consistent Hartree approximation, that the fluctuations destroy the mean-field instability and lead to a first-order phase transition. Theories of nucleation and growth have been constructed based on the idea that the static, Brazovskii-renormalized theory can provide an effective potential for a stochastic Langevin equation. In this paper we examine the validity of this description by using numerical simulations to study the relaxational dynamics of a model described by the Brazovskii Hamiltonian and comparing the results to the predictions of the “static” nucleation theories and to the predictions of a dynamically renormalized theory.

## II The Equation of Motion and Numerical Simulations

Our starting point is the same as in Swift and Hohenberg and describes a system with model A dynamics and a Brazovskii Hamiltonian such as would apply to the diblock copolymers. The Brazovskii Hamiltonian is characterized by a fluctuation spectrum whose maximum occurs at a non-zero wave vector $`|𝐪|=q_0`$ and can be represented by a hypersphere in $`d`$ dimensions. The full form of the Hamiltonian is

$$H=𝑑𝐱[\frac{1}{2}\varphi (^2+q_0^2)^2\varphi +\frac{1}{2}\tau \varphi ^2+\frac{1}{4!}\lambda \varphi ^4]$$
(1)
The dynamics is taken to be relaxational and so the equation of motion is given by the Langevin Equation.

$$\frac{d\varphi }{dt}=M\frac{\delta H}{\delta \varphi }+\eta $$
(2)
where $`M`$ is the mobility (which sets the time scale for the problem) and $`\eta `$ is Gaussian noise ($`<\eta >=0,`$ $`<\eta (\stackrel{}{x},t)`$ $`\eta (\stackrel{}{x}^{},t^{})>=\delta (\stackrel{}{x}\stackrel{}{x}^{})\delta (tt^{})`$ ).

The stochastic Langevin equation derived from the Hamiltonian differs from the usual Ginzburg-Landau description because of the appearance of the unusual gradient term. This dynamical equation is usually referred to as the Swift-Hohenberg equation and falls into the Type $`I_s`$ classification of Cross and Hohenberg . In this classification the system is unstable to a static, spatially periodic structure.

The complete equation of motion is;

$$\frac{d\varphi }{dt}=M(q_0^2^2\varphi +^4\varphi +(q_0^4+\tau )\varphi +\frac{\lambda }{6}\varphi ^3)+\eta $$
(3)
In Eq. (3), the coupling constant, $`\lambda `$ has been rescaled by the noise strength. For systems where the noise strength is small, such as Rayleigh-Benard convection, the effective coupling constant is also small. In di-block copolymers, the noise strength is of the order of $`k_BT`$.

For numerical calculations, the Langevin equation must be approximated by a discrete equation.

$$\varphi (t+\mathrm{\Delta }t)=\varphi (t)M\mathrm{\Delta }t(\mathrm{\pounds }\varphi +\mathrm{\pounds }^2\varphi +(r+q_0^4)\varphi +\frac{\lambda }{6}\varphi ^3)+(M\mathrm{\Delta }t/(\mathrm{\Delta }x)^3)^{1/2}\eta $$
(4)
£ here is the discrete Laplacian. This can take on several forms though for this work we use the simplest form: $`\frac{d^2\varphi }{dx^2}\mathrm{\pounds }\varphi =\frac{1}{\mathrm{\Delta }x^2}(x+\mathrm{\Delta }x)+\varphi (x\mathrm{\Delta }x)2\varphi (x))`$ where the sum is over the dimensions of the lattice. Other choices which include next nearest or more complicated neighborhoods are possible. This is important if isotropy is of concern however the effect is small and can be ignored. The scaling of the noise in Eq. 4 reflects the fact that as the cell size or time steps become smaller (larger), the possible fluctuations become larger (smaller).

For these simulations we set $`q_0=1`$ and choose the lattice spacing such that one lamellar spacing spans six lattice sites ($`\mathrm{\Delta }x=6/2\pi `$) For most of this work the overall lattice size is $`60^3`$. The mobility, $`M`$, is set to one. In order for the simulations to be numerically stable the time scale must be smaller then some stability time, $`\mathrm{\Delta }t1/\mathrm{\Delta }x^4.`$ To satisfy this the timescale is chosen to be $`\mathrm{\Delta }t`$ = 0.01 and measurements are taken at intervals of $`100\mathrm{\Delta }t`$ or longer.

Previous theories of nucleation have analyzed the above equation under the approximation that the effects of fluctuations can be incorporated into a static renormalized free energy $`F_R`$ which replaces the bare free energy $`F`$ in Eq. (2) . We will compare our numerical results to the Brazovskii predictions for the static renormalized parameters and show that the static scenario is beautifully borne out by the simulations in three dimensions. We will then analyze the time-dependence of the structure factor as observed in the numerical simulations and show that this time dependence is inconsistent with the dynamical picture based on a static renormalized free energy. A theory based on a dynamical renormalization of Eq. (2) can provide a qualitatively correct description of the numerical simulations.

## III Statics

### A Brazovskii Theory

Brazovski’s treatment of the model within the Hartree approximation can be restated in terms of an expansion of the thermodynamic potential $`\mathrm{\Gamma }(\overline{\varphi })`$, the generating functional for the vertex functions. This approach has been described in detail by Fredrickson and Binder, and within the Hartree approximation leads to a renormalized mass term ($`r`$) in Eq.(1). The diagrams up to one loop are shown in Fig. 1, and the mass renormalization relation is given by;

$$r+(q^2q_0^2)^2=\tau +(q^2q_0^2)^2+\frac{\lambda }{2}\frac{d𝐪}{(2\pi )^d}(\tau +(q^2q_0^2)^2)^1$$
(5)
This approximation is made self consistent by replacing the bare parameter in the integrand by the renormalized parameter. Essentially the bare propagator on the loop in Fig. (1) is replaced by renormalized propagator. The leads to the Hartree result.

$$r=\tau +\frac{\lambda }{2}\frac{d𝐪}{(2\pi )^d}(r+(q^2q_0^2)^2)^1=\tau +\alpha \lambda /\sqrt{r}$$
(6)
where $`\alpha `$ is proportional to the surface area of a $`d`$ dimensional sphere of unit radius. According to Brazovski this approximation is good only for $`\lambda ^6<<1.`$

The interesting point about Eq. (6) is that even for negative values of the bare parameter, $`\tau `$, the renormalized parameter, $`r`$, is positive. This implies that the disordered phase is always metastable. Brazovskii went on to show that the bare coupling parameter $`\lambda `$ gets renormalized to a negative value leading to a $`\varphi ^6`$ theory and the possibility of a first-order phase transition.

### B Computational Results.

The numerical simulations can measure the static structure factor in the disordered phase which is predicted by the Brazovskii theory  to be,

$$S(q)^1=r+(q^2q_0^2)^2$$
(7)
where $`r`$ is the renormalized control parameter . Thus, $`S(q_0)^1`$ is just $`r`$. The renormalized mass can, therefore, be measured by monitoring the peak of the structure factor. The Hartree calculations predict $`r`$ to be positive and asymptotically approaching zero as $`\tau \mathrm{}`$. Experiments in symmetric di-block copolymers have verified that the behavior of $`S(q_0)`$ is consistent with the Brazovskii predictions and is very different from the mean-field prediction ( $`S(q_0)=\tau ^1`$).

Fig. 2 shows results for $`r`$ obtained from our simulations for different values of the coupling constant, $`\lambda `$. The Gaussian case ($`\lambda =0`$) is also shown for comparison. The system was run for 1000 time steps (each time step being $`100\mathrm{\Delta }t`$) and 100 samples were taken. The data points plotted are the average of the samples while the error bars represent the standard deviation. For large positive values of the bare control parameter, $`\tau `$, the fourth order term becomes less important and $`r`$ approaches the Gaussian value. This is seen clearly for small values of $`\lambda `$. As $`\tau `$ approaches zero, the measured values of $`r`$ deviate from the Gaussian case and stay positive for all $`\tau `$. Some care was taken to normalize the values presented in Fig. 2. For the Gaussian case, $`S(q_0)^1`$ should be linearly related to $`\tau `$ and the slope of the line should be one. In calculating $`S(q_0)`$ several normalizations are needed including the normalization due to the Fourier Transform (FT) and the normalization due to circular averaging. While the FT normalization is just related to the system size, the normalization due to circular averaging is more complicated to calculate. Instead of calculating these normalizations directly, the slope of the raw result from the Gaussian case was used to normalize all of the data. Another concern is that for an infinite continuous system this line should have an x-intercept at zero; however, for the finite systems on a grid used in the simulations the x-intercept is slightly negative. To account for this, all of the data is shifted over by an amount equal to that intercept. The importance of this shift will become apparent in the next paragraph when a scaling is applied and the values of the results around zero are magnified. A small negative (positive) value becomes a much larger negative (positive) value and so shifts from negative to positive will become important.

In reference the authors show that, within the Hartree approximation, the $`r(\tau )`$curves for different $`\lambda ^{}s`$ could be described by a single functional form in terms of scaled variables $`\tau ^{}`$ and $`r^{}`$;

$`r^{}`$ $`=`$ $`r(\alpha \lambda )^{2/3}`$ (8)
$`\tau ^{}`$ $`=`$ $`\tau (\alpha \lambda )^{2/3}`$ (9)
Figure 2 shows plots of $`r^{}`$ against $`\tau ^{}`$obtained from simulations for the three different values of $`\lambda `$. The curves are seen to scale quite well but the scaled curve falls significantly below the theoretical prediction  shown as the thick black line. The theoretical prediction that the value of $`r^{}`$(and $`r`$) never becomes negative and is asymptotic to zero is, however, clearly borne out by the simulations. The deviation from the theoretical line could be a system size effect, however, results for larger and smaller systems are consistent with the data shown in Fig.3. For the Gaussian case $`r^{}=\tau ^{}`$ and for large positive values of the bare parameter both the theory and simulations approach this. Since the theory is larger then the simulations and both are above the Gaussian, our simulations suggest that the contributions to the two point function from diagrams not included in the Hartree approximation serve to lower the overall correction to mean field theory. Another property that the simulation data exhibits in Fig. 3 which is not predicted by the theory is that the curves for different values of $`\lambda `$ diverge from each other at negative values of $`\tau `$, i.e. the scaling is not perfect. Smaller values of $`\lambda `$ lie closer to the theory as expected. The computations were done using values of $`\lambda 10^2`$ and larger while the theory is valid only for $`\lambda 10^6`$ so the theory’s scaling predictions seem to be far more robust then expected.

To further explore the nature of the phase transition in this model, we can compare scaled values of $`S(q_0)`$ obtained from a hot (random) and cold (purely modulated order) initial configurations. Fig. 4 shows just such a plot. As in the previous plot, the theory is shown as a thick black line. For each value of $`\lambda `$ there are two lines presented: one for a disordered start and one for an ordered start. The averaging was done as described above. For positive and small negative values of $`\tau `$ the results from the disordered start and the ordered start are nearly the same. For $`\tau ^{}`$ below some $`\tau _s^{}`$ the two differ with the ordered phase having a higher peak. It should be pointed out that for these values of $`\tau ^{}`$ the disordered state is not in equilibrium (metastable or stable)and there is a slow but definite time evolution of the structure factor over the time period in which the averages are taken. This time-evolution will be discussed in more detail in the section on dynamics. For now the disordered start points are included in the plot simply as a comparison to see where the ordered phase melts. The value where the two data sets diverge can be considered an estimate for the limit of stability (spinodal) of the lamellar phase. As a consistency check the average value of the wave vector was measured. For values of $`\tau `$ above $`\tau _s`$ the wave vector is small and points in a random direction, while for lower values of $`\tau `$ the average wave vector is large and and points along the direction that the system was prepared in. For all three sets of data the lamellar spinodal lies in the range of $`2.1<\tau _s^{}<1.9`$. This is consistent with the prediction from Hartree theory that the first-order transition to the lamellar phase occurs at a lower value of $`\tau ^{}`$, $`\tau _t^{}=2.74`$. That is to say that our measured value of the lamellar spinodal is above the predicted transition temperature as it should be. For values of $`\tau ^{}`$ below $`\tau _s^{}`$, however, systems prepared in the disordered phase always evolve towards the ordered phase. It is possible that the lamellar ordered phase is enhanced by the boundary conditions - the lattice size is always chosen to be commensurate with the lamellar wavelength - however, it should be noted that lamella form in many possible directions where the system size is not commensurate with the lamellar wavelength. Discussion in the dynamics sections should shed some light on this issue.

## IV Dynamics

### A Relaxational Dynamics Based on Static Renormalized Parameters

Previous analysis of nucleation and metastability, in the context of fluctuation-driven phase transitions, have been based on a Langevin equation with the force obtained from the static, Hartree-renormalized free energy function. Fredrickson and Binder used this approach to compute the nucleation barrier and the completion rate of nucleation and growth in di-block copolymers. The interfacial tension was found to be small leading to a small nucleation barrier and rapid nucleation for deep quenches. In this picture, there is no essential difference between the kinetics of a fluctuation-driven first-order transition and a weak-first order transition. The results of our simulations suggest a very different scenario for the growth of the lamellar structures in a Brazovskii model. The shape of the nucleating droplet was more carefully analyzed by Swift and Hohenberg by taking into account spatial inhomogeneities in the effective free energy function. Their analysis relied on constructing a coarse-grained free energy functional. The coarse graining was based on a momentum-shell renormalization idea where fluctuations with momenta far away from the shell defined by $`|𝐪|=𝐪_\mathrm{𝟎}`$ are successively integrated out. This analysis led to non-spherical droplets and were consistent with the picture of spinodal nucleation that had been obtained earlier.

As pointed out in the work of Swift and Hohenberg, a complete theory of the dynamics of fluctuation-driven first order transitions would have to be based on a coarse-graining of the full dynamics as expressed by the original Langevin equation with the bare Brazovskii Hamiltonian. In this paper, we compare the results of our numerical simulations to the predictions of a Hartree-renormalization of the full dynamical equation as expressed in Eq. (2). Our emphasis has been on understanding the nature of the dynamics of the metastable phase and we have not analyzed, in any detail, the spatial structures associated with the growth process.

### B Dynamical Renormalization and Simulation Results

In order to systematically analyze the effects of fluctuations on the kinetics of growth of the lamellar phase, perturbative techniques analogous to the static Hartree approximation have to be applied to the Langevin equation. This is most conveniently done through the dynamical-action formalism. The application of this method to the S-H equation is outlined by Ignatiev et. al. . In the dynamical-action formalism, the average value of an operator $`O(\varphi )`$ over the noise history is rewritten as a functional integral;

$$O\{\varphi (𝐱,t)\}=D\varphi \mathrm{exp}(S[\varphi ]),$$
(10)
where $`S[\varphi ]`$ is the dynamical action.

As in the static case, the calculation of dynamical correlation functions is most conveniently formulated through the construction of a generating function. To establish the closest analogy with the static calculations, it is useful to work with the generator of vertex functions, $`\mathrm{\Gamma }_{dyn}[\overline{\varphi }]`$, and establish a diagrammatic expansion which is the exact analog of the diagrams retained by Brazovskii in the static calculation. The simulation results demonstrate that the static properties of the S-H equation closely follow the Brazovskii predictions and, therefore, it seems appropriate to apply this approximation scheme to the dynamics. The correspondence between statics and dynamics becomes particularly transparent in a super-field formulation which shows that the dynamical perturbation theory in terms of the super fields has exactly the same structure as the static perturbation theory except for the appearance of a very different “kinetic” term which leads to a bare propagator that is distinct from the static theory. The super-field correlation functions can be calculated by constructing diagrams as in the static theory but replacing the static bare propagator by the appropriate dynamical one. The super-field correlation functions encode all the dynamical correlations of the field $`\varphi `$ and the latter can be extracted from well-defined relations. The details of this calculation and the results will be described in a separate publication. In this paper we present the main results which can be compared to the numerical simulations.

The two correlation functions which appear in the dynamical study are

$`G(t,t^{})=<\varphi (t)\eta (t^{})>`$

and

$`C(t,t^{})=<\varphi (t)\varphi (t^{})>`$

Both of these can be obtained from a single super-field correlation function ($`Q(1,2)`$) which is related to $`\mathrm{\Gamma }_{dyn}`$ through

$$Q^1(1,2)=\frac{\delta ^2\mathrm{\Gamma }_{dyn}}{\delta \mathrm{\Phi }_1\delta \mathrm{\Phi }_2}$$
(11)
where $`\mathrm{\Phi }`$’s are the super fields and the derivative is evaluated with the source term set to zero. Applying this formalism to Eq. (2), leads to an expression for $`\mathrm{\Gamma }_{dyn}`$ which involves two-point correlation functions of the fluctuations and is a natural extension of the static $`\mathrm{\Gamma }(\overline{\varphi })`$ to correlation functions involving time. The expression, however, does not lend itself to easy analysis except in two cases, early times when a variant of linear theory can be applied and late times when the system is stationary. The late time dynamics involves the analysis of the nonlinear terms and will be discussed in the concluding section. The early time dynamics is where one is justified in retaining only quadratic terms in the renormalized $`\mathrm{\Gamma }_{dyn}`$. In this limit we obtain the following equation for the equal-time correlation function $`C_q(t,t)=<\varphi _q(t)\varphi _q(t)>`$ which is the structure factor that is monitored in the simulations:

$`C_q(t,t)`$ $`=`$ $`{\displaystyle _0^t}[G(t,t^{})]^2𝑑t^{}`$ (12)
$`G_q(t,t^{})`$ $`=`$ $`{\displaystyle _t^t^{}}(r(t^{\prime \prime })+(q^2q_0^2)^2)𝑑t^{\prime \prime }.`$ (13)
The mass parameter, $`r(t)`$, which is now time-dependent, is renormalized in the dynamic theory using the same approximation as in the static theory. The diagrams used are given in fig. 7. These again are only to one loop. The mass term then becomes

$`r(t)+(q^2q_0^2)^2`$ $`=`$ $`\tau +(q^2q_0^2)^2+`$ (15)
$`{\displaystyle \frac{\lambda }{2}}[{\displaystyle \frac{d𝐪}{(2\pi )^d}\frac{1}{(\tau +(q^2q_0^2)^2)}}{\displaystyle \frac{d𝐪}{(2\pi )^d}\frac{\mathrm{exp}\{2(\tau +(q^2q_0^2)^2)t\}}{(\tau +(q^2q_0^2)^2)}}]`$
Again, the approximation is made self consistent by replacing the $`\tau `$ with $`r`$ in the integrand.

$$r(t)=r_{eq}\frac{\lambda }{2}\frac{d^3q}{(2\pi )^3}\frac{\mathrm{exp}(2D(q)t)}{D(q)}$$
(16)
$`D(q)`$ here is the renormalized propagator, $`D(q)=r_{eq}+(q^2q_0^2)^2`$ where $`r_{eq}`$ denotes the static renormalized mass parameter which is the solution to eq. 7.

In Eq. 16, at long times the integrand in the second term becomes small and $`r(t)`$ approaches $`r_{eq}.`$ which is confirmed in the section on statics above. At time zero the subtraction of the second term is just a subtraction of the Hartree level correction factor to the bare parameter introduced earlier Eq. (6 and, therefore, $`r(t=0)`$ is just the bare parameter $`\tau `$, while as $`t\mathrm{}`$, $`r(t)`$ approaches the renormalized, equilibrium value.

This scenario is confirmed by the simulations. The simulation results discussed in Section III B confirm that the equilibrium value, $`r_{eq}`$, is consistent with the Hartree prediction. Analyzing the early stage dynamics can verify the value of $`r(t)`$ at short times. Figure 5 shows the growth of $`S(q_0)`$ as a function of time for various values of the bare control parameter, $`\tau `$. To average over noise, five independent runs where taken for each value of the control parameter. If $`r(t)`$ is time-independent, the early-stage evolution (linear theory) is described by

$$S(q,t)=C_q(t,t)=S(q,0)\mathrm{exp}\{D(q)t\}+\frac{1}{D(q)}(1\mathrm{exp}\{D(q)t\})$$
(17)
where $`D(q)`$ is $`D(q)=r_0+(q^2q_0^2)^2`$ and $`r_0`$ is the time-independent value of $`r`$. At the peak of the structure factor, $`q=q_0`$ , and $`D(q_0)`$ is just $`r_0`$, the value of which can be estimated by fitting the simulation results shown in Fig. 5 to 17. The results of such fits are presented in Fig. 6. This figure shows that $`r_0`$ and $`\tau `$ are linearly related with a slope which is nearly one. This implies that for short times, the dynamically renormalized parameter is linearly related to the bare parameter and is not renormalized to positive values for negative values of the bare parameter. The very early-stage dynamics, therefore, is characteristic of a system exhibiting unstable growth.

These fits must be considered with some care. In the case where $`r`$ is not time dependent, the linear theory describes the system only for short times, that is times on the order of the natural time scale, $`r^1`$ \[See for example. For the fastest growing mode, when the data is fit to different time scales the results obtained are consistent for all time scales up to $`r^1`$. For the S-H equation, however, the results obtained are dependent on the time range. In particular, as the range of time increases the values obtained for $`r`$ become less negative. This is consistent with $`r(t)`$ growing as a function of time and the results of the fits merely represent some average value of $`r(t)`$ for the time range involved. If the system has been quenched to a negative value of $`\tau `$, Eq. 16 indicates that the early-time evolution will exhibit unstable growth. As time evolves the second term in Eq. (16) decreases and the value of $`r(t)`$ approaches $`r_{eq}`$ which is always positive. Though the integral in Eq. 16 is hard to evaluate analytically, numerics can provide some insight into its behavior. Figure 8 shows a numerical evaluation of the time evolution of $`r(t)`$ for various values of $`r_{eq}`$ and $`\lambda =0.01`$. As confirmed in the statics section above, $`r(t)`$ will eventually become positive no matter how deep the quench is, though the time for this to happen may become very long. As long as $`r(t)`$ is negative the system is expected to undergo unstable growth though, since the mass parameter is time-dependent, the time evolution will be more complicated than the usual spinodal decomposition scenario. We can define a crossover time, $`t_{cross}`$ at which $`r(t)`$ changes sign, and an average growth time, $`t_{ave}`$. The average growth time is just the inverse of $`r_{ave}`$ defined as

$$t_{ave}^1=r_{ave}=\frac{1}{t_{cross}}_0^{t_{cross}}r(t^{})𝑑t^{}$$
(18)
In Fig. 9, we show plots of $`t_{cross}`$ and $`t_{ave}`$ as a function of the bare control parameter. For shallow quenches $`t_{cross}`$ is small compared to $`t_{ave}`$ and the system reaches a well-defined metastable equilibrium state which is disordered. If $`t_{cross}`$ becomes larger then $`t_{ave}`$, metastability becomes difficult to define. The system takes longer to equilibrate in the metastable disordered phase than it takes to grow lamellar structures. We can use the condition $`t_{cross}=t_{ave}`$ to define a crossover temperature $`\tau _{dyn}^{}`$. Above this temperature the disordered phase quickly becomes locally stable and a nucleation event is needed for the formation of the lamellar structures. Below $`\tau _{dyn}^{}`$, the system is expected to evolve continuously towards the lamellar phase without any evidence of metastability of the disordered phase. Fig. 9 is in sharp contrast to standard Ginzburg-Landau (G-L) theory. For the time dependent G-L equation, the presence of noise suppresses the critical point. In the region where the bare parameter is negative but the renormalized parameter is positive, $`t_{cross}`$ is always smaller then $`t_{ave}`$ and so there is no unstable growth until the true critical point is reached.

For the system under consideration here, different scenarios are possible depending on the relationship of this dynamical crossover temperature to the first-order transition temperature, $`\tau _{trans}^{}`$ , at which the free energy of the lamellar phase becomes lower that that of the disordered phase. If $`\tau _{trans}^{}`$ is higher than $`\tau _{dyn}^{}`$, there will be a regime of temperatures over which the the system will undergo nucleation and growth. On the other hand if $`\tau _{trans}^{}`$ is lower than $`\tau _{dyn}^{}`$, then there is no nucleating regime and one would observe only unstable growth, albeit of an unusual nature since the free-energy surface is evolving with time. From figure 9 we can see that $`t_{cross}`$ is small only for very shallow quenches below the mean-field transition. The results of our simulations, presented in the next section, are in qualitative agreement with this scenario. We would like to emphasize that the growth dynamics described above are qualitatively different from rapid nucleation in which the system quickly reaches the stable equilibrium state. A large value of $`t_{cross}`$, on the other hand, indicates a type of ergodicity breaking as the system takes a very long time to reach equilibrium.

All of our numerical results indicate that $`\tau _{trans}^{}`$ predicted from Hartree is lower than $`\tau _{dyn}^{}`$ and they are remarkably close to each other. We have been unable to come up with a simple relationship between these two temperatures and, therefore, can only interpret the similarity of the two as a remarkable coincidence. The dynamical crossover is deduced from time scales which characterize the evolution of the free-energy surface while the transition temperature is deduced from a comparison of the depth of the two wells. It is not clear why the two temperatures should be similar in magnitude. It can be argued that there is only a single parameter, $`\lambda `$, controlling the scale of fluctuations and, therefore, the two temperatures should be related, however, there is no obvious argument to suggest that they should be identical.

The picture emerging from the dynamical renormalization is a natural extension of the effect of fluctuations on the static results of the Brazovskii model. Following a quench from a relatively high temperature to a temperature where $`\tau `$ is negative, the system is in a locally unstable region (top of a hill). As the fluctuations grow with time, the non-linear terms characterized by $`\lambda `$ become important and they renormalize the curvature of the hill. This scenario is quite different from the usual picture of evolution in an adiabatic potential.

### C Simulation Results for Late Times

Further evidence for the dynamical scenario presented in the previous section comes from examining the long time evolution of the peak of the structure factor obtained from our numerical simulations and correlating that to snap shots of the system as it evolves. Figure 10 is a plot of the amplitude of the structure factor peak as it evolves. For shallow quenches, $`\tau >0.04,`$the peak grows to an equilibrium value and does not evolve further. These values are consistent with the equilibrium values reported above in the section on statics. For deeper quenches the peak amplitude grows quickly and then appears to saturate; however careful examination shows that the value continues to grow very slowly. This is consistent with late stage domain growth which as been studied extensively by Elder et. al. for 2-D systems . The last feature in the system is a final rise to an equilibrium value. This rise is a finite size affect is caused by the majority domain finally taking over the entire system. This being the case, the peak would not be expected to grow much since the difference between the value before and after this final evolution should only represent the surface area between the different domains.

For a quench depth of $`\tau =0.07`$ this time evolution can be compared to a series of system snapshots shown in figures 11\- 15. These figures represent $`\varphi =0`$ isosurfaces which would be the boundaries between the different microphases of the system. For early times the system appears to be very disordered while at later times domains of ordered lamellar structures begin to appear. Just before the final convulsive growth, at a time $`t=4000`$, there still appear to be a few different ordered domains while just after the final growth, at $`t=5000`$, only one domain appears to be left in the system. Thus we still have a consistent picture of slow but continuous domain growth in the system which is governed by the evolution of the 2nd order term as it goes from negative values to positive values.

## V Future Work

In the work presented above we have shown that analysis of the static free energy does not always provide an adequate description of the fluctuation driven dynamics. Though figure 10 shows features which are reminiscent of a first order phase transition, dynamical analysis of the renormalized coefficients suggest that a more complicated evolution is taking place. A similar situation occurs for the superconducting transition which is also a fluctuation-driven first-order transition.

In the present work we have paid close attention to the time evolution of the mass parameter and its effect on the evolution of the system at early times. The dynamical renormalization results also become tractable at times large compared to $`t_{cross}`$ when the system has reached equilibrium. As pointed out earlier, one enters this regime fairly quickly for shallow quenches. At this stage, the higher order terms in $`\mathrm{\Gamma }_{dyn}`$ become important. The most interesting aspect of the renormalized theory is that the fourth-order term becomes non-local in time :

$`\mathrm{\Gamma }_{dyn}(4){\displaystyle d1d2\mathrm{\Phi }^2(1)K(12)\mathrm{\Phi }^2(2)}`$

and the integral of the kernel over all times leads to the Brazovskii result for the renormalized fourth-order term. One can derive an “effective” Langevin equation for the fields $`\varphi `$ from the renormalized $`\mathrm{\Gamma }_{dyn}`$, and the above equation implies that this Langevin equation has “memory” effects which appear in the non-linear term. This would lead to unusual behavior of the equilibrium correlations functions and might provide a distinctive signature of fluctuation-driven first-order transitions. We are in the process of investigating the effective equation for $`\varphi `$ numerically. One obvious consequence of the non-local term is that the barrier to nucleation is dynamic and it is not valid to think of nucleation as tunneling through a fixed barrier. The situation is closer to the problem of quantum tunneling for many degrees of freedom where one also finds an effective equation which is non-local in time.

Another area of interest is the role of defects in the phase transition of the system. These will be important for very shallow quenches to just below the transition point and we have noted there presence when looking at snapshots of the system. These defects are reminiscent of the perforated lamellar phase described in reference . One way to study this is to use a directional order parameter measure introduced by Bray et al. to study the 2D problem. This order parameter is similar to the order parameter for complex fluids and provides directional as well as density information. What it may show is that the local gradient terms are enhanced for negative values $`\tau `$ suggesting that perhaps there is a disordered to nematic transition in this region which others have suggested for 2D .

We would like to close with a discussion of possible experimental systems where this dynamical scenario can be studied experimentally. Since the effective $`\lambda `$ is extremely small in Rayleigh-Benard systems the regime where the disordered state is metastable will probably be unaccessible. In di-block copolymers, however, there should be a significant temperature range, below the mean-field transition, where the disordered state is locally stable and a study of the equilibrium, two-time correlation function should reveal the existence of the non-linear memory term. It might also be possible to observe unusual nucleation if the system parameters allow for $`\tau _{trans}^{}`$ to be higher than $`\tau _{dyn}^{}`$.