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# Roth’s Theorem in the Primes

## 1. Introduction

Arguably the second most famous result of Klaus Roth is his 1953 upper bound  on $`r_3(N)`$, defined 17 years previously by Erdős and Turán to be the density of the largest set $`A[N]`$ containing no non-trivial 3-term arithmetic progression (3AP). Roth was the first person to show that $`r_3(N)=o(1)`$. In fact, he proved the following quantitative version of this statement.

###### Proposition 1.1 (Roth).

We have $`r_3(N)1/\mathrm{log}\mathrm{log}N`$.

There was no improvement on this bound for nearly 40 years, until Heath-Brown  and Szemerédi  proved that $`r_3(\mathrm{log}N)^c`$ for some small positive constant $`c`$. Recently Bourgain  provided the best bound currently known.

###### Proposition 1.2 (Bourgain).

We have $`r_3(N)\left(\mathrm{log}\mathrm{log}N/\mathrm{log}N\right)^{1/2}`$.

The methods of Heath-Brown, Szemerédi and Bourgain may be regarded as (highly non-trivial) refinements of Roth’s technique. There is a feeling that Proposition 1.2 is close to the natural limit of this method. This is irritating, because the sequence of primes is not covered by these results. However it is known that the primes contain infinitely many 3APs.111In April 2004 the author and T. Tao published a preprint shoing that the primes contain arbitrarily long arithmetic progressions.

###### Proposition 1.3 (Van der Corput).

The primes contain infinitely many *3AP*s.

Van der Corput’s method is very similar to that used by Vinogradov to show that every large odd number is the sum of three primes. Let us also mention a paper of Balog  in which it is shown that for any $`n`$ there are $`n`$ primes $`p_1,\mathrm{},p_n`$ such that all of the averages $`\frac{1}{2}(p_i+p_j)`$ are prime. In this paper we propose to prove a common generalization of the results of Roth and Van der Corput. Write $`𝒫`$ for the set of primes.

###### Theorem 1.4.

Every subset of $`𝒫`$ of positive upper density contains a *3AP*.

In fact, we get an explicit upper bound on the density of a 3AP-free subset of the primes, but it is ridiculously weak. Observe that as an immediate consequence of Theorem 1.4 we obtain what might be termed a van der Waerden theorem in the primes, at least for progressions of length 3. That is, if one colours the primes using finitely many colours then one may find a monochromatic 3AP.

We have not found a written reference for the question answered by Theorem 1.4, but M.N. Huxley has discussed it with several people .

To prove Theorem 1.4 we will use a variant of the following result. This says that the primes enjoy what is known as the Hardy-Littlewood majorant property.

###### Theorem 1.5.

Suppose that $`p2`$ is a real number, and let $`𝒫_N=𝒫[1,N]`$. Let $`\{a_n\}_{n𝒫_N}`$ be any sequence of complex numbers with $`|a_n|1`$ for all $`n`$. Then

$$\underset{n𝒫_N}{}a_ne(n\theta )_{L^p(𝕋)}C(p)\underset{n𝒫_N}{}e(n\theta )_{L^p(𝕋)},$$
(1.1)
where the constant $`C(p)`$ depends only on $`p`$.

It is perhaps surprising to learn that such a property does not hold with any set $`\mathrm{\Lambda }[N]`$ in place of $`𝒫_N`$. Indeed, when $`p`$ is an even integer it is rather straightforward to check that any set does satisfy (1.1) (with $`C(p)=1`$). However, there are sets for which (1.1) fails badly when $`p`$ is not an even integer. For a discussion of this see  and for related matters including connections with the Kakeya problem, see .

We will apply a variant of Theorem 1.5 for $`p=5/2`$, when it certainly does not seem to be trivial. To prove it, we will establish a somewhat stronger result which we call a restriction theorem for primes. The reason for this is that our argument is very closely analogous to an argument of Tomas and Stein  concerning Fourier transforms of measures supported on spheres.

A proof of the restriction theorem for primes was described, in a different context, by Bourgain . Our argument, being visibly analagous to the approach of Tomas, is different and has more in common with §3 of . This more recent paper of Bourgain deals with restriction phenomena of certain sets of lattice points.

To deduce Theorem 1.4 from (a variant of) Theorem 1.5 we use a variant of the technique of granularization as developed by I.Z. Ruzsa and the author in a series of papers beginning with , as well as a “statistical” version of Roth’s theorem due to Varnavides. We will also require an argument of Marcinkiewicz and Zygmund which allows us to pass from the continuous setting in results such as (1.1) – that is to say, $`𝕋`$ – to the discrete, namely $`/N`$.

Finally, we would like to remark that it is possible, indeed probable, that Roth’s theorem in the primes is true on grounds of density alone. The best known lower bound on $`r_3(N)`$ comes from a result of Behrend  from 1946.

###### Proposition 1.6 (Behrend).

We have $`r_3(N)e^{C\sqrt{\mathrm{log}N}}`$ for some absolute constant $`C`$.

This may well give the correct order of magnitude for $`r_3(N)`$, and if anything like this could be proved Theorem 1.4 would of course follow trivially.

## 2. Preliminaries and an outline of the argument

Although the main results of this paper concern the primes in $`[N]`$, it turns out to be necessary to consider slightly more general sets. Let $`m\mathrm{log}N`$ be a positive integer and let $`b`$, $`0bm1`$, be coprime to $`m`$. We may then define a set

$$\mathrm{\Lambda }_{b,m,N}=\left\{nN\right|nm+b\text{is prime}\}.$$

We expect $`\mathrm{\Lambda }_{b,m,N}`$ to have size about $`mN/\varphi (m)\mathrm{log}N`$, and so it is natural to define a function $`\lambda _{b,m,N}`$ supported on $`\mathrm{\Lambda }_{b,m,N}`$ by setting

$$\lambda _{b,m,N}(n)=\{\begin{array}{cc}\varphi (m)\mathrm{log}(nm+b)/mN\hfill & \text{if }n\mathrm{\Lambda }_{b,m,N}\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$

For simplicity we write $`X=\mathrm{\Lambda }_{b,m,N}`$ for the next few pages. We will abuse notation and consider $`\lambda _{b,m,N}`$ as a measure on $`X`$. Thus for example $`\lambda _{b,m,N}(X)`$, which is defined to be $`_n\lambda _{b,m,N}(n)`$, is roughly $`1`$ by the prime number theorem in arithmetic progressions. We use $`L^p(d\lambda _{b,m,N})`$ norms and also the inner product $`f,g_X=f(n)\overline{g(n)}\lambda _{b,m,N}(n)`$ without further comment.

It is convenient to use the wedge symbol for the Fourier transforms on both $`𝕋`$ and $``$, which we define by $`f^{}(n)=f(\theta )e(n\theta )𝑑\theta `$ and $`g^{}(\theta )=_ng(n)e(n\theta )`$ respectively. Here, of course, $`e(\alpha )=e^{2\pi i\alpha }`$.

For any measure space $`Y`$ let $`B(Y)`$ denote the space of continuous functions on $`Y`$ and define a map $`T:B(X)B(𝕋)`$ via

$$T:f(f\lambda _{b,m,N})^{}.$$
(2.1)
The object of this section is to give a new proof of the following result, which may be a called a restriction theorem for primes.

###### Theorem 2.1 (Bourgain).

Suppose that $`p>2`$ is a real number. Then there is a constant $`C(p)`$ such that for all functions $`f:X`$ we have

$$Tf_pC(p)N^{1/p}f_2.$$
(2.2)
Remember that the $`L^2`$ norm is taken with respect to the measure $`\lambda _{b,m,N}`$. Theorem 2.1 probably has most appeal when $`b=m=1`$, in which case we may derive consequences for the primes themselves. Later on, however, we will take $`m`$ to be a product of small primes, and so it is necessary to have the more general form of the theorem.

We turn now to an outline of the proof of Theorem 2.1. The analogy between our proof and an argument by Tomas , giving results of a similar nature for spheres in high-dimensional Euclidean spaces, is rather striking. In fact, the reader may care to look at the presentation of Tomas’s proof in , whereupon she will see that there is an almost exact correspondence between the two arguments.

To begin with, the proof proceeds by the method of $`T`$ and $`T^{}`$, a basic technique in functional analysis. One can check that the operator $`T^{}:B(𝕋)B(X)`$ is given by

$$T^{}:gg^{}|_X,$$
(2.3)
by verifying the relation

$$Tf,g_𝕋=(f\lambda _{b,m,N})^{}(\theta )\overline{g(\theta )}𝑑\theta =\underset{n}{}f(n)\overline{g^{}(n)}\lambda _{b,m,N}(n)=f,T^{}g_X.$$

The equation (2.3) explains the term restriction. Using (2.3) we see that the operator $`TT^{}`$ is the map from $`B(𝕋)`$ to itself given by

$$TT^{}:ff\lambda _{b,m,N}^{}.$$
(2.4)
Now Theorem 2.1 may be written, in an obvious notation, as

$$T_{2p}C(p)N^{1/p}.$$
(2.5)
The principle of $`T`$ and $`T^{}`$, as we will use it, states that

$$T_{2p}^2=TT^{}_{p^{}p}=T^{}_{p^{}2}^2$$
(2.6)
We would like to emphasise that there is nothing mysterious going on here – this result is just an elegant and convenient way of bundling together some applications of Hölder’s inequality. The proof of the part that we will need, that is to say is the inequality $`T_{2p}^2TT^{}_{p^{}p}`$, is simply

$`Tf_p`$ $`=`$ $`\underset{g_p^{}=1}{sup}Tf,g`$
$`=`$ $`\underset{g_p^{}=1}{sup}f,T^{}g`$
$``$ $`f_2\underset{g_p^{}=1}{sup}T^{}g_2`$
$`=`$ $`f_2\underset{g_p^{}=1}{sup}g,TT^{}g^{1/2}`$
$``$ $`f_2TT^{}_{p^{}p}^{1/2}.`$
Thus we will, for much of the paper, be concerned with showing that the operator $`TT^{}`$ as given by (2.4) satisfies the bound

$$TT^{}_{p^{}p}C^{}(p)N^{2/p}.$$
(2.7)
The preceding remarks show that a proof of this will imply Theorem 2.1. To get such a bound one splits $`\lambda `$ into certain dyadic pieces, that is a sum

$$\lambda _{b,m,N}=\underset{j=1}{\overset{K}{}}\psi _j+\psi _{K+1}.$$
(2.8)
The slightly curious way of writing this indicates that the definition of $`\psi _{K+1}`$ will be a little different from that of the other $`\psi _j`$. We will define these pieces so that they satisfy the $`L^1`$$`L^{\mathrm{}}`$ estimates

$$f\psi _j^{}_{\mathrm{}}_ϵ2^{(1ϵ)j}f_1$$
(2.9)
for some $`ϵ<(p2)/2`$, and also the $`L^2`$$`L^2`$ estimates

$$f\psi _j^{}_2_ϵ\frac{2^{ϵj}}{N}f_2.$$
(2.10)
Applying the Riesz-Thorin interpolation theorem (see , Chapter 7) will then give

$$f\psi _j^{}_p2^{\delta j}N^{2/p}f_p^{}$$

for some positive $`\delta `$ (depending on $`ϵ`$). Summing these estimates from $`j=1`$ to $`K+1`$ will establish (2.7) and hence Theorem 2.1.

To define the decomposition (2.8) we need yet more notation. From the outset we will suppose that we are trying to prove Theorem 2.1 for a particular value of $`p`$ – the argument is highly and essentially non-uniform in $`p`$. Write $`A=4/(p2)`$. Let $`1<Q(\mathrm{log}N)^A`$. If $`b,m,N`$ are as before (recall that $`m\mathrm{log}N`$) then we define a measure $`\lambda _{b,m,N}^{(Q)}`$ on $``$ by setting

$$\lambda _{b,m,N}^{(Q)}(n)=\{\begin{array}{cc}N^1_{\begin{array}{c}pQ\\ pm\end{array}}\left(1\frac{1}{p}\right)^1\hfill & \text{if }nN\text{ and }p|(nm+b)p>Q\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$

Define $`\lambda _{b,m,N}^{(1)}(n)=0`$ for all $`n`$.

As $`Q`$ becomes large the measures $`\lambda _{b,m,N}^{(Q)}`$ look more and more like $`\lambda _{b,m,N}`$. Much of §4 will be devoted to making this principle precise. We will sometimes refer to the support of $`\lambda _{b,m,N}^{(Q)}`$ as the set of $`Q`$-rough numbers.

Now let $`K`$ be the smallest integer with

$$2^K>\frac{1}{10}(\mathrm{log}N)^A$$
(2.11)
and define

$$\psi _j=\lambda _{b,m,N}^{(2^j)}\lambda _{b,m,N}^{(2^{j1})}$$
(2.12)
for $`j=1,\mathrm{},K`$ and define

$$\psi _{K+1}=\lambda _{b,m,N}\lambda _{b,m,N}^{(2^K)},$$
(2.13)
so that (2.8) holds. In the next two sections we prove the two required estimates, (2.9) and (2.10).

Let us note here that the main novelty in our proof of Theorem 2.1 lies in the definition of the dyadic decomposition (2.8). By contrast, the analogous dyadic decompositions in  take place on the Fourier side, requiring the introduction of various smooth cutoff functions not specifically related to the underlying arithmetic structure.

## 3. An $`L^2`$$`L^2`$ estimate

It turns out that the proof of (2.10), the $`L^2`$$`L^2`$ estimate, is by far the easier of the two estimates required. We have

$`f\psi _j^{}_2`$ $`=`$ $`\widehat{f}\psi _j_2`$
$``$ $`\psi _j_{\mathrm{}}\widehat{f}_2`$
$`=`$ $`\psi _j_{\mathrm{}}f_2.`$
Suppose first of all that $`1jK`$. Then

$`\psi _j_{\mathrm{}}`$ $``$ $`\lambda _{b,m,N}^{(2^j)}_{\mathrm{}}+\lambda _{b,m,N}^{(2^{j1})}_{\mathrm{}}`$
$`=`$ $`N^1{\displaystyle \underset{\begin{array}{c}p2^{j+1}\\ pm\end{array}}{}}\left(1{\displaystyle \frac{1}{p}}\right)^1+N^1{\displaystyle \underset{\begin{array}{c}p2^j\\ pm\end{array}}{}}\left(1{\displaystyle \frac{1}{p}}\right)^1.`$
The two products here may be estimated using Merten’s formula \[14, Chapter 22\]:

$$\underset{pQ}{}(1p^1)\frac{e^\gamma }{\mathrm{log}Q}.$$

This gives

$$\psi _j_{\mathrm{}}j/N,$$
(3.1)
and hence

$$f\psi _j^{}_2\frac{j}{N}f_2,$$
(3.2)
which is certainly of the requisite form (2.10). For $`j=K+1`$ we have

$`\psi _{K+1}_{\mathrm{}}`$ $``$ $`\lambda _{b,m,N}^{(2^K)}_{\mathrm{}}+\lambda _{b,m,N}_{\mathrm{}}`$
$``$ $`\mathrm{log}N/N,`$
so that

$$f\psi _{K+1}^{}_2\frac{\mathrm{log}N}{N}f_2.$$
(3.3)
This also constitutes an estimate of the type (2.10) for some $`ϵ<(p2)/2`$. Indeed, recalling our choice of $`A`$ and $`K`$ (viz. (2.11)) one can check that $`2^K(\mathrm{log}N)^{1/ϵ}`$ for some such $`ϵ`$.

## 4. An $`L^1`$$`L^{\mathrm{}}`$ estimate

This section is devoted to the rather lengthy task of proving estimates of the form (2.9).

Introduction. The first step towards obtaining an estimate of the form (2.9) is to observe that

$$f\psi _j^{}_{\mathrm{}}\psi _j^{}_{\mathrm{}}f_1.$$
(4.1)
We will prove that $`\psi _j^{}_{\mathrm{}}`$ is not too large by proving

###### Proposition 4.1.

Suppose that $`Q(\mathrm{log}N)^A`$. Then we have the estimate

$$\lambda _{b,m,N}^{}\lambda _{b,m,N}^{(Q)}_{\mathrm{}}\mathrm{log}\mathrm{log}Q/Q.$$

The detailed proof of this fact will occupy us for several pages. Let us begin, however, by using (4.1) to see how it implies an estimate of the form (2.9). If $`1jK`$ then we have

$`\psi _j^{}_{\mathrm{}}`$ $`=`$ $`\lambda _{b,m,N}^{(2^j)}\lambda _{b,m,N}^{(2^{j1})}_{\mathrm{}}`$ (4.2)
$``$ $`\lambda _{b,m,N}^{}\lambda _{b,m,N}^{(2^j)}_{\mathrm{}}+\lambda _{b,m,N}^{}\lambda _{b,m,N}^{(2^{j1})}_{\mathrm{}}`$
$``$ $`\mathrm{log}j/2^j.`$
This is certainly of the form (2.9). The estimate for $`j=K+1`$ is even easier, being immediate from Proposition 4.1.

To prove Proposition 4.1 we will use the Hardy-Littlewood circle method. Thus we divide $`𝕋`$ into two sets, traditionally referred to as the major and minor arcs. It is perhaps best if we define these explicitly at the outset. Thus let $`p`$ be the exponent for which we are trying to prove Theorem 2.1. Recall that $`A=4/(p2)`$, and set $`B=2A+20`$. These numbers will be fixed throughout the proof. By Dirichlet’s theorem on approximation, every $`\theta 𝕋`$ satisfies

$$\left|\theta \frac{a}{q}\right|\frac{(\mathrm{log}N)^B}{qN}$$
(4.3)
for some $`qN(\mathrm{log}N)^B`$ and some $`a`$, $`(a,q)=1`$. The major arcs consist of those $`\theta `$ for which $`q`$ can be taken to be at most $`(\mathrm{log}N)^B`$. We will write this collection using the notation

$$𝔐=\underset{\begin{array}{c}q(\mathrm{log}N)^B\\ (a,q)=1\end{array}}{}𝔐_{a,q}.$$

For these $`\theta `$, the Fourier transforms $`\lambda _{b,m,N}^{(Q)}`$ and $`\lambda _{b,m,N}^{}`$ depend on the distribution of the almost-primes and primes along arithmetic progressions with common difference at most $`(\mathrm{log}N)^B`$. The minor arcs $`𝔪`$ consist of all other $`\theta `$. Here different techniques apply, and one can conclude that both $`\lambda _{b,m,N}^{(Q)}`$ and $`\lambda _{b,m,N}^{}`$ are small. The triangle inequality then applies.

The ingredients are as follows. The almost-primes are eminently suited to applications of sieve techniques. To keep the paper as self-contained as possible, we will follow Gowers  and use arguably the simplest sieve, that due to Brun, on both the major and minor arcs.

The genuine primes, on the other hand, are harder to deal with. Here we will quote two well-known results from the literature. The information concerning distribution along arithmetic progressions to small moduli comes from the prime number theorem of Siegel and Walfisz.

###### Proposition 4.2 (Siegel–Walfisz).

Suppose that $`q(\mathrm{log}N)^B`$, that $`(a,q)=1`$ and that $`1N_1N_2N`$. Then

$$\underset{\begin{array}{c}N_1<pN_2\\ pa(\text{mod}q)\end{array}}{}\mathrm{log}p=\frac{N_2N_1}{\varphi (q)}+O\left(N\mathrm{exp}(C_B\sqrt{\mathrm{log}N})\right).$$

The rather strange formulation of the theorem reflects the fact that the constant $`C_B`$ is ineffective for any $`B1`$ due to the possible existence of a Siegel zero. For more information, including a complete proof of Proposition 4.2, see Davenport’s book .

The techniques for dealing with the minor arcs are associated with the names of Weyl, Vinogradov and Vaughan.

The major arcs. We will have various functons $`f:[N]`$ with

$$f_{\mathrm{}}=O(\mathrm{log}N/N)$$
(4.4)
which are regularly distributed along arithmetic progressions in the following sense. If $`LN(\mathrm{log}N)^{2BA1}`$ and if $`X[N]`$ is an arithmetic progression $`\{r,r+q,\mathrm{},r+(L1)q\}`$ with $`q(\mathrm{log}N)^B`$ then

$$\underset{nX}{}f(n)=\frac{L}{N}\left(\gamma _{r,q}(f)+O((\mathrm{log}N)^A)\right),$$
(4.5)
where $`\gamma _{r,q}`$ depends only on $`r`$ and $`q`$, $`|\gamma _{r,q}|q`$ and the implied constant in the $`O`$ term is absolute. This information is enough to get asymptotics for $`f^{}(\theta )`$ when $`|\theta a/q|`$ is small, as we prove in the next few lemmas.

For a residue $`r`$ modulo $`q`$, write $`N_r`$ for the set $`\{nN:nr(\text{mod}q)\}`$. Write $`\tau `$ for the function on $`𝕋`$ defined by $`\tau (\theta )=N^1_{nN}e(\theta n)`$. The first lemma deals with $`f^{}(\theta )`$ for $`|\theta |(\mathrm{log}N)^B/qN`$.

###### Lemma 4.3.

Let $`r`$ be a residue modulo $`q`$, suppose that $`|\theta |(\mathrm{log}N)^B/qN`$, and suppose that the function $`f`$ satisfies (4.4) and (4.5). Then

$$\underset{nN_r}{}f(n)e(\theta n)=q^1\gamma _{r,q}(f)\tau (\theta )+O(q^1(\mathrm{log}N)^A).$$

Proof. Set $`L=N(\mathrm{log}N)^{2BA1}`$ and partition $`N_r`$ into arithmetic progressions $`(X_i)_{i=1}^T`$ of common difference $`q`$ and length between $`L`$ and $`2L`$, where $`T2N/Lq`$. For each $`i`$ fix an element $`x_iX_i`$.

$`{\displaystyle \underset{nN_r}{}}f(n)e(\theta n)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{T}{}}}{\displaystyle \underset{nX_i}{}}f(n)e(\theta n)`$ (4.6)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{T}{}}}e(\theta x_i){\displaystyle \underset{nX_i}{}}f(n)+{\displaystyle \underset{i=1}{\overset{T}{}}}{\displaystyle \underset{nX_i}{}}f(n)\left(e(\theta n)e(\theta x_i)\right)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{T}{}}}e(\theta x_i){\displaystyle \frac{|X_i|}{N}}\left(\gamma _{r,q}(f)+O((\mathrm{log}N)^A)\right)+O(LN^1q^1(\mathrm{log}N)^{B+1})`$
$`=`$ $`\gamma _{r,q}(f){\displaystyle \underset{i=1}{\overset{T}{}}}e(\theta x_i){\displaystyle \frac{|X_i|}{N}}+O\left(q^1(\mathrm{log}N)^A\right).`$
However

$`{\displaystyle \underset{i=1}{\overset{T}{}}}e(\theta x_i)|X_i|`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{T}{}}}{\displaystyle \underset{nX_i}{}}e(\theta n)+{\displaystyle \underset{i=1}{\overset{T}{}}}{\displaystyle \underset{nX_i}{}}\left(e(\theta x_i)e(\theta n)\right)`$ (4.7)
$`=`$ $`{\displaystyle \underset{nN_r}{}}e(n\theta )+O(Lq^1(\mathrm{log}N)^B).`$
Finally, observe that if $`0r,sq1`$ then

$$\underset{nN_r}{}e(\theta n)\underset{nN_s}{}e(\theta n)=O((\mathrm{log}N)^B),$$

and so

$$\left|N^1\underset{nN_r}{}e(\theta n)q^1\tau (\theta )\right|=O(N^1(\mathrm{log}N)^B).$$

Combining this with (4.6) and (4.7) completes the proof of the lemma.

We may now get an asymptotic for $`f^{}(\theta )`$ when $`\theta `$ is in the neighbourhood of $`a/q`$.

###### Lemma 4.4.

Suppose that $`f`$ satisfies the conditions (4.4) and (4.5) and that $`\theta 𝔐_{a,q}`$ for some $`a,q`$ with $`(a,q)=1`$ and $`q(\mathrm{log}N)^B`$. Write

$$\sigma _{a,q}(f)=\underset{r}{}e(ar/q)\gamma _{r,q}(f).$$
(4.8)
Then we have

$$f^{}(\theta )=q^1\sigma _{a,q}(f)\tau (\theta a/q)+O((\mathrm{log}N)^A).$$
(4.9)
Proof. Write $`\beta =\theta a/q`$. Then

$`f^{}(\theta )`$ $`=`$ $`{\displaystyle \underset{nN}{}}f(n)e(\theta n)`$
$`=`$ $`{\displaystyle \underset{r(\text{mod}q)}{}}e(ar/q){\displaystyle \underset{nN_r}{}}f(n)e(\beta n)`$
$`=`$ $`q^1\tau (\beta ){\displaystyle \underset{r(\text{mod}q)}{}}e(ar/q)\gamma _{r,q}(f)+O((\mathrm{log}N)^A)`$
$`=`$ $`q^1\sigma _{a,q}(f)\tau (\beta )+O((\mathrm{log}N)^A).`$
This concludes the proof of the lemma.

To apply these lemmas, we need to show that $`f=\lambda _{b,m,N}^{(Q)}`$ and $`f=\lambda _{b,m,N}`$ satisfy (4.4) and (4.5) for suitable choices of $`\gamma _{r,q}(f)`$. We will then evaluate the sums $`\sigma _{a,q}(f)`$. This slightly tedious business is the subject of our next four lemmas.

###### Lemma 4.5.

$`f=\lambda _{b,m,N}`$ satisfies (4.4) and (4.5) with

$$\gamma _{r,q}(f)=\{\begin{array}{cc}\varphi (m)q/\varphi (mq)\hfill & \text{if }(mr+b,mq)=1\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$

Proof. This is a fairly immediate consequence of the Siegel-Walfisz Theorem (Proposition 4.2). Let $`X=\{r,r+q,\mathrm{},r+(L1)q\}`$ be any progression contained in $`[N]`$ with common difference $`q(\mathrm{log}N)^B`$ and length $`LN(\mathrm{log}N)^{2BA1}`$. An element $`r+jqX`$ lies in $`\mathrm{\Lambda }_{b,m,N}`$ precisely if $`(mr+b)+jmq`$ is prime, so the lemma is trivially true unless $`(mr+b,mq)=1`$. Supposing this to be the case, we may use Proposition 4.2. Recalling that $`m\mathrm{log}N`$, one has

$`\lambda _{b,m,N}(X)`$ $`=`$ $`{\displaystyle \frac{\varphi (m)qL}{\varphi (mq)N}}+O\left(mq\mathrm{exp}(C_{B+1}\sqrt{\mathrm{log}mqN})\right)`$
$`=`$ $`{\displaystyle \frac{L}{N}}\left({\displaystyle \frac{\varphi (m)q}{\varphi (mq)}}+O((\mathrm{log}N)^A)\right),`$
as required.

###### Lemma 4.6.

$`f=\lambda _{b,m,N}^{(Q)}`$ satisfies (4.4) and (4.5) with

$$\gamma _{r,q}(f)=\{\begin{array}{cc}_{\begin{array}{c}pQ\\ pm\end{array}}\left(1\frac{1}{p}\right)^1_{\begin{array}{c}pQ\\ pmq\end{array}}\left(1\frac{1}{p}\right)\hfill & \text{if }(mr+b,mq)\text{ is }Q\text{-rough}\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$

Proof. Consider an arithmetic progression $`X=\{r,r+q,\mathrm{},r+(L1)q\}`$. Let $`p_1,\mathrm{},p_k`$ be the primes with $`pQ`$ and $`pm`$. If $`(mr+b,mq)`$ is not $`Q`$-rough then $`p_i|(mr+b,mq)`$ for some $`i`$, and the second alternative of the lemma clearly holds. Suppose then that $`(mr+b,mq)`$ is $`Q`$-rough. We will apply the Brun sieve to estimate $`\lambda _{b,m,N}^{(Q)}(X)`$.

Let $`xX`$ be chosen uniformly at random, and for each $`i`$ let $`X_i`$ be the event $`p_i|(mx+b)`$. Since $`p_i(mr+b,mq)`$, the probability of $`X_i`$ is $`ϵ_i/p_i+O(L^1)`$, where $`ϵ_i=0`$ if $`p_i|q`$ and $`ϵ_i=1`$ otherwise. Now we have

$$\frac{N}{L}\underset{\begin{array}{c}pQ\\ pm\end{array}}{}\left(1\frac{1}{p}\right)\lambda _{b,m,N}^{(Q)}(X)=\left(X_i^c\right)=U,$$
(4.10)
say. By the inclusion-exclusion formula it follows that for every positive integer $`t`$

$$U=\underset{s=0}{\overset{t}{}}(1)^s\underset{1i_1<\mathrm{}<i_sk}{}\underset{j=1}{\overset{s}{}}ϵ_{i_j}/p_{i_j}+O(L^1)\underset{s=1}{\overset{t}{}}\left(\genfrac{}{}{0pt}{}{k}{s}\right).$$
(4.11)
It is helpful to have the error term here in a more usable form. To this end, observe that it is certainly at most $`O(k^t/L)`$. We wish to replace the main term in (4.11) by $`_{i=1}^k\left(1ϵ_i/p_i\right)`$, which is equal to the completed sum

$$\underset{s=0}{\overset{k}{}}(1)^s\underset{1i_1<\mathrm{}<i_sk}{}\underset{j=1}{\overset{s}{}}ϵ_{i_j}/p_{i_j}.$$

Doing this introduces an error

$$E=\underset{s=t+1}{\overset{k}{}}(1)^s\underset{1i_1<\mathrm{}<i_sk}{}\underset{j=1}{\overset{s}{}}ϵ_{i_j}/p_{i_j},$$

which is bounded above by

$$\underset{s=t+1}{\overset{k}{}}\frac{1}{s!}\left(\underset{i=1}{\overset{k}{}}\frac{1}{p_i}\right)^s.$$
(4.12)
By another result of Mertens one has $`_{i=1}^kp_i^1\mathrm{log}\mathrm{log}Q+O(1)`$. Hence if $`t3\mathrm{log}\mathrm{log}Q`$ then each term in (4.12) is at most one half the previous one, leading to the bound

$$|E|\frac{2(\mathrm{log}\mathrm{log}Q)^t}{t!}\left(\frac{4e\mathrm{log}\mathrm{log}Q}{t}\right)^t.$$

Combining all of this gives

$$U=\underset{i=1}{\overset{k}{}}\left(1ϵ_i/p_i\right)+O(k^t/L)+O\left((4e\mathrm{log}\mathrm{log}Q/t)^t\right).$$

Using the trivial bound $`kQ`$, and choosing $`t=\mathrm{log}N/2A\mathrm{log}\mathrm{log}N`$, one gets

$`U`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}\left(1ϵ_i/p_i\right)+O(N^{1/4A})`$
$`=`$ $`{\displaystyle \underset{\begin{array}{c}pQ\\ pmq\end{array}}{}}\left(1{\displaystyle \frac{1}{p}}\right)+O(N^{1/4A}).`$
The lemma is immediate from this and (4.10); we have

$`\lambda _{b,m,N}^{(Q)}(X)`$ $`=`$ $`{\displaystyle \underset{\begin{array}{c}pQ\\ pm\end{array}}{}}\left(1{\displaystyle \frac{1}{p}}\right)^1{\displaystyle \frac{L}{N}}\left({\displaystyle \underset{\begin{array}{c}pQ\\ pmq\end{array}}{}}\left(1{\displaystyle \frac{1}{p}}\right)+O(N^{1/4A})\right)`$
$`=`$ $`{\displaystyle \frac{L}{N}}\left(\gamma _{r,q}+O((\mathrm{log}N)^A)\right),`$
where $`\gamma _{r,q}`$ has the form claimed.

Building on the last lemma, the next lemma gives an evaluation of $`\sigma _{a,q}(\lambda _{b,m,N}^{(Q)})`$ and an asymptotic for $`\lambda _{b,m,N}^{(Q)}(\theta )`$ when $`\theta 𝔐_{a,q}`$. If $`Q2`$ we say that a positive integer is $`Q`$-smooth if all of its prime divisors are at most $`Q`$. We declare there to be no $`1`$-smooth numbers.

###### Lemma 4.7.

Suppose that $`(a,q)=1`$. Then

$$\sigma _{a,q}(\lambda _{b,m,N}^{(Q)})=\{\begin{array}{cc}\frac{q\mu (q)}{\varphi (q)}e\left(\frac{ab\overline{m}}{q}\right)\hfill & \text{if }(m,q)=1\text{ and }q\text{ is }Q\text{-smooth;}\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array}$$

where $`\overline{m}`$ is the inverse of $`m`$ modulo $`q`$. If $`\theta 𝔐_{a,q}`$ then

$$\lambda _{b,m,N}^{(Q)}(\theta )=\{\begin{array}{cc}\frac{\mu (q)}{\varphi (q)}e\left(\frac{ab\overline{m}}{q}\right)\tau \left(\theta \frac{a}{q}\right)+O((\mathrm{log}N)^A)\hfill & \text{if }(m,q)=1\hfill \\ & \text{and }q\text{ is }Q\text{-smooth;}\hfill \\ O\left((\mathrm{log}N)^A\right)\hfill & \text{otherwise}.\hfill \end{array}$$

Proof. Recall the definition (4.8) of $`\sigma _{a,q}`$, and also Lemma 4.6. We shall prove that

$$\underset{\begin{array}{c}r(\text{mod}q)\\ (mr+b,mq)\text{is }Q\text{-rough}\end{array}}{}e(ar/q)=\{\begin{array}{cc}e(ab\overline{m}/q)\mu (q)\hfill & \text{if }(m,q)=1\text{ and }q\text{ is }Q\text{-smooth}\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(4.13)
Now if $`p|m`$ then $`p`$ can never divide $`mr+b`$, because we are assuming that $`(m,b)=1`$. Let $`q_0`$ be the largest factor of $`q`$ which is a product of primes $`p`$ with $`pQ`$ and $`pm`$. Then the sum (4.13) is just

$$\underset{\begin{array}{c}r(\text{mod}q)\\ (q_0,mr+b)=1\end{array}}{}e(ar/q).$$
(4.14)
Set $`q_1=q/q_0`$ and write, for each $`r`$ mod $`q`$, $`r=kq_0+s`$ where $`0kq_11`$ and $`s`$ is a residue mod $`q_0`$. Then the sum (4.14) is

$$\underset{\begin{array}{c}s(\text{mod}q_0)\\ (q_0,ms+b)=1\end{array}}{}\underset{k=0}{\overset{q_11}{}}e\left(\frac{a(kq_0+s)}{q}\right)=\underset{\begin{array}{c}s(\text{mod}q_0)\\ (q_0,mr+b)=1\end{array}}{}e(as/q)\underset{k=0}{\overset{q_11}{}}e(ak/q_1).$$

Now $`a`$ is coprime to $`q`$ and hence to $`q_1`$, and therefore the rightmost sum here vanishes unless $`q_1=1`$. This is the case precisely if $`q_0=q`$, which means that $`(q,m)=1`$ and $`q`$ is $`Q`$-smooth. In this case, the sum is

$$\underset{\begin{array}{c}s(\text{mod}q)\\ (q,ms+b)=1\end{array}}{}e(as/q).$$

Set $`t=ms+b`$. Then this sum is just

$`{\displaystyle \underset{\begin{array}{c}t(\text{mod}q)\\ (q,t)=1\end{array}}{}}e\left({\displaystyle \frac{a\overline{m}(tb)}{q}}\right)`$ $`=`$ $`e(ab\overline{m}/q){\displaystyle \underset{(q,t)=1}{}}e(a\overline{m}t/q)`$
$`=`$ $`e(ab\overline{m}/q)\mu (q).`$
This last evaluation, of what is known as a Ramanujan Sum, is well-known and is contained, for example, in . This proves (4.13).

Now to obtain $`\sigma _{a,q}`$ we must simply multiply (4.13) by the factor

$$F=\underset{\begin{array}{c}p<Q\\ pm\end{array}}{}\left(1\frac{1}{p}\right)^1\underset{\begin{array}{c}pQ\\ pmq\end{array}}{}\left(1\frac{1}{p}\right)$$

appearing in Lemma 4.6. One gets zero unless $`(m,q)=1`$ and $`q`$ is $`Q`$-smooth, in which case it is not hard to see that $`F=q/\varphi (q)`$. This completes the evaluation of $`\sigma _{a,q}(\lambda _{b,m,N}^{(Q)})`$, and the claimed form for $`\lambda _{b,m,N}^{(Q)}(\theta )`$ is an immediate consequence of Lemma 4.4.

We need a version of the above lemma in which $`\lambda _{b,m,N}^{(Q)}`$ is replaced by $`\lambda _{b,m,N}`$. Fortunately, we can save ourselves some work by noticing that for fixed $`q`$ and $`m`$ we have

$$\gamma _{r,q}(\lambda _{b,m,N})=\gamma _{r,q}(\lambda _{b,m,N}^{(Q)})$$
(4.15)
for sufficiently222Here we regard $`\gamma _{r,q}(\lambda _{b,m,N})`$ and $`\gamma _{r,q}(\lambda _{b,m,N}^{(Q)})`$ as purely formal expressions, so there is no issue of whether or not, for example, Lemma 4.7 is valid for “sufficiently large” $`Q`$. large $`Q`$. Thus $`\sigma _{a,q}(\lambda _{b,m,N})`$ can be evaluated by simply letting $`Q\mathrm{}`$ in the first formula of Lemma 4.7. We get

$$\sigma _{a,q}(\lambda _{b,m,N})=\{\begin{array}{cc}q\mu (q)e(ab\overline{m}/q)\hfill & \text{if }(q,m)=1\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(4.16)
This immediately leads, via Lemma 4.4, to the following evaluation of $`\lambda _{b,m,N}^{}(\theta )`$.

###### Lemma 4.8.

Suppose that $`(a,q)=1`$ and that $`\theta 𝔐_{a,q}`$. then

$$\lambda _{b,m,N}^{}(\theta )=\{\begin{array}{cc}\frac{\mu (q)}{\varphi (q)}e\left(\frac{ab\overline{m}}{q}\right)\tau \left(\theta \frac{a}{q}\right)+O\left((\mathrm{log}N)^A\right)\hfill & \text{if }(m,q)=1\hfill \\ O\left((\mathrm{log}N)^A\right)\hfill & \text{otherwise}.\hfill \end{array}$$
(4.17)
The minor arcs. In this subsection we look at $`\lambda _{b,m,N}^{}(\theta )`$ and $`\lambda _{b,m,N}^{(Q)}(\theta )`$ when $`\theta `$ is not close to a rational with small denominator.

###### Lemma 4.9.

Suppose that $`a,q`$ are positive integers with $`(a,q)=1`$, and let $`\theta `$ be a real number such that $`|\theta a/q|1/q^2`$. Then

$$\lambda _{b,m,N}^{}(\theta )(\mathrm{log}N)^{10}\left(q^{1/2}+N^{1/5}+N^{1/2}q^{1/2}\right).$$
(4.18)
Thus if $`\theta 𝔪`$ then $`\lambda _{b,m,N}^{}(\theta )=O((\mathrm{log}N)^A)`$.

Remarks. This is a well-known estimate, at least when $`b=m=1`$. The first (unconditional) results of this type were obtained by I.M. Vinogradov, and nowadays it is possible to give a rather clean argument thanks to the identity of Vaughan . Chapter 24 of Davenport’s book  describes the use of Vaughan’s identity in the more general context of the estimation of sums $`_{nN}\mathrm{\Lambda }(n)f(n)`$. To obtain Lemma 4.9 we used this approach, but could afford to obtain results which are rather non-uniform in $`m`$ due to the restriction $`m\mathrm{log}N`$ under which we are operating. Details may be found in the supplementary document . We remark that existing results in the literature concerning minor arcs estimates for primes restricted to arithmetic progressions, such as , strive for a much better dependence on the parameter $`m`$.

###### Lemma 4.10.

Suppose that $`a,q`$ are positive integers with $`(a,q)=1`$, and let $`\theta `$ be a real number such that $`|\theta a/q|1/q^2`$. Then

$$\lambda _{b,m,N}^{(Q)}(\theta )(\mathrm{log}N)^3\left(q^1+qN^1+N^{1/8A}\right).$$
(4.19)
Thus if $`\theta 𝔪`$ then $`\lambda _{b,m,N}^{(Q)}(\theta )=O((\mathrm{log}N)^A)`$.

Proof. Let $`p_1,\mathrm{},p_k`$ be the primes less than or equal to $`Q`$ which do not divide $`m`$. Another application of the inclusion-exclusion principle gives

$$\lambda _{b,m,N}^{(Q)}(\theta )=N^1e(b\theta /m)\underset{i=1}{\overset{k}{}}\left(1\frac{1}{p_i}\right)^1h(\theta ),$$

where

$$h(\theta )=\underset{s=0}{\overset{k}{}}(1)^s\underset{1i_1<\mathrm{}<i_sk}{}\underset{\begin{array}{c}1yNm/p_{i_1}\mathrm{}p_{i_s}\\ yb(\text{mod}m)\end{array}}{}e\left(\frac{\theta p_{i_1}\mathrm{}p_{i_s}y}{m}\right).$$
(4.20)
Summing the geometric progression, one sees that the inner sum is no more than

$$\mathrm{min}\{\theta p_{i_1}\mathrm{}p_{i_s}^1,2mN/p_{i_1}\mathrm{}p_{i_s}\}.$$

We will split the sum over $`s`$ in (4.20) into two pieces, over the ranges $`s[0,t]`$ and $`s(t,k]`$ where $`t=\mathrm{log}N/2A\mathrm{log}\mathrm{log}N`$. Each of the primes $`p_i`$ is at most $`Q(\mathrm{log}N)^A`$, so the product of any $`st`$ of them is no more than $`\sqrt{N}`$. Of course, all such products are distinct and so

$$\underset{s=0}{\overset{t}{}}(1)^s\underset{1i_1<\mathrm{}<i_sk}{}\underset{\begin{array}{c}yNm/p_{i_1}\mathrm{}p_{i_s}\\ yb(\text{mod}m)\end{array}}{}e\left(\frac{\theta p_{i_1}\mathrm{}p_{i_s}y}{m}\right)\underset{n\sqrt{N}}{}\mathrm{min}(\theta n^1,2mN/n).$$

This is a quantity whose estimation is standard in this area because of its pertinence to the estimation of exponential sums on minor arcs. It is bounded above by $`C(\mathrm{log}N)^3(N^{1/2}+q+Nq^1)`$; details may once again be found in .

On the other hand

$`{\displaystyle \underset{s=t+1}{\overset{k}{}}}(1)^s{\displaystyle \underset{1i_1<\mathrm{}<i_sk}{}}{\displaystyle \underset{\begin{array}{c}yNm/p_{i_1}\mathrm{}p_{i_s}\\ yb(\text{mod}m)\end{array}}{}}e\left({\displaystyle \frac{\theta p_{i_1}\mathrm{}p_{i_s}y}{m}}\right)`$
$``$ $`2mN{\displaystyle \underset{s=t+1}{\overset{k}{}}}{\displaystyle \underset{1i_1<\mathrm{}<i_sk}{}}{\displaystyle \underset{j=1}{\overset{s}{}}}p_{i_j}^1`$
$``$ $`2mN{\displaystyle \underset{s=t+1}{\overset{k}{}}}(s!)^1\left(p_1^1+\mathrm{}+p_k^1\right)^s`$
$``$ $`4mN(2e\mathrm{log}\mathrm{log}\mathrm{log}N/t)^tmN^{11/4A}N^{11/8A}.`$
Since $`_{i=1}^k(11/p_i)^1\mathrm{log}N`$, the claimed bound follows.

Proof of Proposition 4.1. Suppose first of all that $`\theta 𝔐_{a,q}`$ for some $`a,q`$, and recall Lemmas 4.7 and 4.8. If $`q`$ is $`Q`$-smooth then

$$\left|\lambda _{b,m,N}^{}(\theta )\lambda _{b,m,N}^{(Q)}(\theta )\right|=O(N(\mathrm{log}N)^A).$$

If $`q`$ is not $`Q`$-smooth then $`q>Q`$ and so we get

$`\left|\lambda _{b,m,N}^{}(\theta )\lambda _{b,m,N}^{(Q)}(\theta )\right|`$ $``$ $`|\lambda _{b,m,N}^{}(\theta )|+|\lambda _{b,m,N}^{(Q)}(\theta )|`$
$``$ $`2/\varphi (q)+O((\mathrm{log}N)^A)`$
$``$ $`4\mathrm{log}\mathrm{log}Q/Q+O((\mathrm{log}N)^A),`$
the last estimate being contained in , Chapter 17. Since we are assuming that $`Q(\mathrm{log}N)^A`$ this expression is $`O(\mathrm{log}\mathrm{log}Q/Q)`$. If, on the other hand, $`\theta 𝔪`$ then we have

$`\left|\lambda _{b,m,N}^{}(\theta )\lambda _{b,m,N}^{(Q)}(\theta )\right|`$ $``$ $`|\lambda _{b,m,N}^{}(\theta )|+|\lambda _{b,m,N}^{(Q)}(\theta )|`$
$`=`$ $`O((\mathrm{log}N)^A)`$
$`=`$ $`O(Q^1).`$
This at last completes the proof of Proposition 4.1.

## 5. Restriction and majorant estimates for primes

In this section we prove Theorems 1.5 and 2.1.

We have already seen, in (4.1) and (4.2), how Proposition 4.1 implies an $`L^1`$$`L^{\mathrm{}}`$ estimate for the operator $`ff\psi _j`$ of the form (2.9). In fact, we have

$$f\psi _j_{\mathrm{}}\frac{\mathrm{log}j}{2^j}f_1$$
(5.1)
for all $`j=1,\mathrm{},K+1`$. For each fixed $`j=1,\mathrm{},K`$, one can use the Riesz-Thorin interpolation theorem to interpolate between (3.2) and (5.1). This theorem, which is discussed in \[11, Chapter 7\], is better known to analytic number theorists as the type of convexity principle that underpins many basic estimates on $`\zeta `$ and $`L`$-functions. It gives

$$f\psi _j_pj^{2/p}(\mathrm{log}j)^{12/p}2^{(12/p)j}N^{2/p}f_p^{}.$$
(5.2)
For $`j=K+1`$ another interpolation, now between (3.3) and (5.1), instead gives

$$f\psi _{K+1}_p(\mathrm{log}N)^{2/p}(\mathrm{log}K)^{12/p}2^{(12/p)K}.$$

Recalling at this point the definition (2.11) of $`K`$ we see that this implies

$$f\psi _{K+1}_p(\mathrm{log}N)^{1/p}N^{2/p}.$$

Summing this together with (5.2) for $`j=1,\mathrm{},K`$ gives, because of the decomposition (2.8),

$$f\lambda _{b,m,N}_pC(p)N^{2/p}f_p^{}.$$

As we have already remarked, Theorem 2.1 follows by the principle of $`T`$ and $`T^{}`$.

Now we prove Theorem 1.5. Although we will need a slightly different result later on, this theorem seems to be the most elegant way to state the majorant property for the primes.

Proof of Theorem 1.5. Let $`(a_n)_{n𝒫_N}`$ be any sequence of complex numbers with $`|a_n|1`$ for all $`n`$. We apply Theorem 2.1 to the function $`f`$ defined by $`f(n)=a_n/\mathrm{log}n`$. Writing out the conclusion of Theorem 2.1 gives, for any $`p>2`$,

$$\left|\underset{n}{}f(n)\mathrm{log}ne(n\theta )\right|^p𝑑\theta _pN^{p/21}\left(\underset{n}{}f(n)^2\mathrm{log}n\right)^{p/2}.$$

Therefore

$`{\displaystyle \left|\underset{n𝒫_N}{}a_ne(n\theta )\right|^p𝑑\theta }`$ $`_p`$ $`N^{p/21}\left({\displaystyle \underset{n𝒫_N}{}}{\displaystyle \frac{|a_n|^2}{\mathrm{log}n}}\right)^{p/2}`$
$`_p`$ $`N^{p1}(\mathrm{log}N)^p.`$
However it is an easy matter to check that

$$\left|\underset{n𝒫_N}{}e(n\theta )\right|^p𝑑\theta _{|\theta |1/2N}\left|\underset{n𝒫_N}{}e(n\theta )\right|^p𝑑\theta N^{p1}(\mathrm{log}N)^p.$$

This proves Theorem 1.5 for $`p>2`$. For $`p=2`$ it is trivial using Parseval’s identity.

## 6. Roth’s theorem in the primes

Let $`A_0`$ be a subset of the primes with positive relative upper density. By this we mean that there is a positive constant $`\alpha _0`$ such that, for infinitely many integers $`n`$, we have

$$|A𝒫_n|\alpha _0n/\mathrm{log}n.$$
(6.1)
This is not a particularly convenient statement to work with, and our first lemma derives something more useful from it.

###### Lemma 6.1.

Suppose that there is a set $`A_0𝒫`$ with positive relative density, but which contains no *3AP*s. Then there is a positive real number $`\alpha `$ and infinitely many primes $`N`$ for which the following is true. There is a set $`A\{1,\mathrm{},N/2\}`$, and an integer $`W[\frac{1}{8}\mathrm{log}\mathrm{log}N,\frac{1}{4}\mathrm{log}\mathrm{log}N]`$ such that

* $`A`$ contains no *3AP*s
* $`\lambda _{b,m,N}(A)\alpha `$ for some $`b`$ with $`(b,m)=1`$, where $`m=_{pW}p`$.

Proof. Take any $`n\alpha _0^3`$ for which (6.1) holds. Let $`W=\frac{1}{4}\mathrm{log}\mathrm{log}n`$, and set $`m=_{pW}p`$. Choose $`N`$ to be any prime in the range $`(2n/m,4n/m]`$. Now there are certainly no more than $`m`$ elements of $`A_0`$ which share a factor with $`m`$, and no more than $`n^{3/4}`$ elements $`xA_0`$ with $`xn^{3/4}`$. Thus

$$\underset{b:(b,m)=1}{}\underset{\begin{array}{c}xn\\ xb(\text{mod}m)\end{array}}{}A_0(x)\mathrm{log}x\alpha _0n/2,$$

and for some choice of $`b`$ we have

$$\underset{\begin{array}{c}xn\\ xb(\text{mod}m)\end{array}}{}A_0(x)\mathrm{log}x\alpha _0n/2\varphi (m).$$
(6.2)
Write $`A=m^1\left((A_0[n])b\right)`$. This set, being a part of $`A_0`$ subjected to a linear transformation, contains no $`3`$-term AP. It is also clear that $`A\{1,\mathrm{},N/2\}`$. Furthermore (6.2) is equivalent to

$$\underset{\begin{array}{c}xN\\ mx+b\text{is prime}\end{array}}{}A(x)\mathrm{log}(mx+b)\alpha _0n/2\varphi (m),$$

which implies that $`\lambda _{b,m,N}(A)\alpha _0n/2mN\alpha _0/8`$. The lemma follows, with $`\alpha =\alpha _0/8`$.  The reason we stipulate that $`A`$ be contained in $`\{1,\mathrm{},N/2\}`$ is that $`A`$ does not contain any 3APs when considered as a subset of $`_N=/N`$. This allows us to make us of Fourier analysis on $`_N`$. If $`f:_N`$ is a function we will write, for any $`r_N`$,

$$\stackrel{~}{f}(r)=\underset{x_N}{}f(x)e(rx/N).$$

Observe that $`f`$ may also be considered as a function on $``$ via the embedding $`_N[N]`$, and then $`\stackrel{~}{f}(r)=f^{}(r/N)`$.

For notational simplicity write $`\mu =\lambda _{b,m,N}`$. We will consider $`A`$ and $`\mu `$ as functions on $`_N`$. Write $`a=A\mu `$. We will continue to abuse notation by using $`\mu `$ and $`a`$ as measures. Thus, for example, $`a(_N)\alpha `$.

Now if $`A`$ contains no (non-trivial) 3APs then

$`{\displaystyle \underset{x,d}{}}a(x)a(x+d)a(x+2d)`$ $`=`$ $`{\displaystyle \underset{x}{}}a(x)^3`$ (6.3)
$``$ $`{\displaystyle \underset{x}{}}\mu (x)^3`$
$``$ $`(\mathrm{log}N)^3/N^2.`$
We are going to show that this forces $`\alpha `$ to be small. We will do this by constructing a new measure $`a_1`$ on $`_N`$ which is set-like, which means that $`a_1`$ behaves a bit like $`N^1`$ times the characteristic function of a set of size $`\alpha N`$. The new measure $`a_1`$ will be fairly closely related to $`a`$, and in fact we will be able to show that

$$\underset{x,d}{}a_1(x)a_1(x+d)a_1(x+2d)\text{is small}.$$
(6.4)
This, it turns out, is impossible; an argument of Varnavides based on Roth’s theorem tells us that a dense subset of $`_N`$ contains lots of 3APs. We will adapt his argument in a trivial way to show that the same is true of set-like measures.

The arguments of this section, then, fall into two parts. First of all we must define $`a_1`$, define the notion of “set-like” and then show that $`a_1`$ is indeed set-like. The key ingredient here is Lemma 6.2, which says that $`\stackrel{~}{\mu }`$ is small away from zero. Secondly, we must formulate and prove a result of the form (6.4). For this we need Theorem 2.1, the restriction theorem for primes.

The idea of constructing $`a_1`$, and the technique for constructing it, has its origins in the notions of granularization as used in a paper of I.Z. Ruzsa and the author . In the present context things look rather different however and, in the absence of anything which might be called a “grain”, we think the terminology of  no longer appropriate.

Let us proceed to the definition of $`a_1`$. Let $`\delta (0,1)`$ be a real number to be chosen later, and set

$$R=\{r_N:|\stackrel{~}{a}(r)|\delta \}.$$

Let $`k=|R|`$, and write $`R=\{r_1,\mathrm{},r_k\}`$. Let $`ϵ(0,1)`$ be another real number to be chosen later, and write $`B(R,ϵ)`$ for the Bohr neighbourhood

$$\{x_N:\frac{xr_i}{N}ϵi[k]\}.$$

Write $`B=B(R,ϵ)`$ and set $`\beta (x)=B(x)/|B|`$. Define

$$a_1=a\beta \beta .$$
(6.5)
It is easy to see that

$$a_1(_N)\alpha .$$
(6.6)
In Lemma 6.3 below we will show that $`a_1_{\mathrm{}}2/N`$, provided that a certain inequality between $`ϵ,k`$ and $`W`$ is satisfied. This is what we mean by the statement that $`a_1`$ is set-like.

###### Lemma 6.2.

Suppose that $`N`$, and hence $`W`$, is sufficiently large. We have

$$\underset{r0}{sup}|\stackrel{~}{\mu }(r)|2\mathrm{log}\mathrm{log}W/W.$$

Proof. Recall that $`\stackrel{~}{\mu }(r)=\mu ^{}(r/N)`$. There are three different cases to consider.

Case 1. $`r/N𝔐_{0,1}`$, that is to say $`|r/N|(\mathrm{log}N)^B/N`$. Then by Lemma 4.8 we have the asymptotic

$$\stackrel{~}{\mu }(r)=\tau (r/N)+O(\mathrm{log}N)^A.$$

Observe, however, that $`\tau (r/N)=0`$ provided that $`r0`$.

Case 2. $`r/N𝔐_{a,q}`$. Then Lemma 4.8 gives

$$\stackrel{~}{\mu }(r)=\frac{\chi _q\mu (q)}{\varphi (q)}e\left(\frac{ab\overline{m}}{q}\right)\tau \left(\frac{r}{N}\frac{a}{q}\right)+O(\mathrm{log}N)^A,$$

where

$$\chi _q=\{\begin{array}{cc}1\hfill & (q,m)=1\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$

Since $`m=_{pW}p`$, we certainly have $`\chi _q=0`$ for $`qW`$. Thus indeed

$$|\stackrel{~}{\mu }(r)|\underset{nW}{sup}\varphi (n)^1+O(\mathrm{log}N)^A2\mathrm{log}\mathrm{log}W/W.$$

Case 3. $`r/N𝔪`$. Then Lemma 4.9 gives $`\stackrel{~}{\mu }(r)=\mu ^{}(r/N)=O((\mathrm{log}N)^A)`$.

###### Lemma 6.3.

Suppose that $`ϵ^k2\mathrm{log}\mathrm{log}W/W`$. Then the measure $`a_1`$ is set-like, in the sense that we have $`a_1_{\mathrm{}}2/N`$.

Proof. Indeed

$`a_1(x)`$ $`=`$ $`a\beta \beta (x)`$
$``$ $`\mu \beta \beta (x)`$
$`=`$ $`N^1{\displaystyle \underset{r}{}}\stackrel{~}{\mu }(r)\stackrel{~}{\beta }(r)^2e(rx/N)`$
$``$ $`N^1\stackrel{~}{\mu }(0)\stackrel{~}{\beta }(0)^2+N^1{\displaystyle \underset{r0}{}}|\stackrel{~}{\mu }(r)||\stackrel{~}{\beta }(r)|^2`$
$``$ $`N^1+N^1\underset{r0}{sup}|\stackrel{~}{\mu }(r)|{\displaystyle \underset{r}{}}|\stackrel{~}{\beta }(r)|^2`$
$`=`$ $`N^1+|B|^1\underset{r0}{sup}|\stackrel{~}{\mu }(r)|`$
$``$ $`N^1+{\displaystyle \frac{2\mathrm{log}\mathrm{log}W}{W|B|}}.`$
Now by a well-known application of the pigeonhole principle we have $`|B|ϵ^kN`$, from which the lemma follows immediately.

We move on now to the second part of our programme, which is a statement and proof of a result of the form (6.4).

###### Proposition 6.4.

We have

$$\underset{x,d}{}a_1(x)a_1(x+d)a_1(x+2d)C^{}N^{3/2}+\frac{1}{N}\left(2^{12}ϵ^2\delta ^{5/2}+C\delta ^{1/2}\right).$$

We will require several lemmas. The most important is a “discrete majorant property”. Before we state and prove this, we give an elegant argument of Marcinkiewicz and Zygmund . We outline the argument here since we like it and, possibly, it is not particularly well-known.

###### Lemma 6.5 (Marcinkiewicz–Zygmund).

Let $`N`$ be a positive integer, and let $`f:[N]`$ be any function. Consider $`f`$ also as a function on $`_N`$. Let $`p>1`$ be a real number. Then we have

$$\underset{r_N}{}|\stackrel{~}{f}(r)|^p=\underset{r=0}{\overset{N1}{}}|f^{}(r/N)|^pC(p)N|\stackrel{~}{f}(\theta )|^p𝑑\theta .$$

Proof. Consider the function

$$g(n)=2\left(1\frac{|n|}{2N}\right)\chi _{|n|2N}\left(1\frac{|n|}{N}\right)\chi _{|n|N}.$$

This function is equal to $`1`$ for all $`n`$ with $`|n|N`$. Its Fourier transform, $`g^{}(\theta )`$, is equal to $`2K_{2N}(\theta )K_N(\theta )`$, a difference of two Fejér kernels. Thus we have

$$f^{}=f^{}\left(2K_{2N}K_N\right),$$

and so

$`|\stackrel{~}{f}(r)|^p`$ $`=`$ $`|f^{}(r/N)|^p`$
$`=`$ $`\left|{\displaystyle f^{}(\theta )\left(2K_{2N}(r/N\theta )K_N(r/N\theta )\right)𝑑\theta }\right|^p`$
$``$ $`3^{p1}\left(2^p\left|{\displaystyle f^{}(\theta )K_{2N}(r/N\theta )𝑑\theta }\right|^p+\left|{\displaystyle f^{}(\theta )K_N(r/N\theta )𝑑\theta }\right|^p\right)`$
$``$ $`3^{p1}\left(2^p{\displaystyle |f^{}(\theta )|^pK_{2N}(r/N\theta )𝑑\theta }+{\displaystyle |f^{}(\theta )|^pK_N(r/N\theta )𝑑\theta }\right)`$
by two applications of Jensen’s inequality. It is necessary, of course, to use the fact that the Fejér kernels are non-negative. To conclude the proof, one only has to show that

$$\underset{r=0}{\overset{N1}{}}K_N(r/N\theta )CN,$$

together with a similar inequality for $`K_{2N}`$. But this is a straightforward matter using the bound

$$\underset{r=0}{\overset{N1}{}}K_N(r/N\theta )\underset{j=0}{\overset{N1}{}}\underset{\varphi [\frac{j}{N},\frac{j+1}{N}]}{sup}K_N(\varphi )$$

together with the estimate

$$K_N(\varphi )\mathrm{min}(N,N^1|\varphi |^2),$$

valid for $`|\varphi |1/2`$.

###### Lemma 6.6 (Discrete majorant property).

Suppose that $`p>2`$. Then there is an absolute constant $`C(p)`$ *(*not depending on $`a`$*)* such that

$$\underset{r}{}|\stackrel{~}{a}(r)|^pC(p).$$

Proof. A direct application of Theorem 2.1 gives

$$|a^{}(\theta )|^p𝑑\theta C^{}(p)N^1.$$

The lemma is immediate from this and Lemma 6.5.

###### Lemma 6.7.

Suppose that $`rR`$. Then $`\left|1\stackrel{~}{\beta }(r)^4\stackrel{~}{\beta }(2r)^2\right|2^{12}ϵ^2`$.

Proof. We have

$`\left|1\stackrel{~}{\beta }(r)\right|`$ $`=`$ $`{\displaystyle \frac{1}{|B|}}\left|{\displaystyle \underset{xB}{}}\left(1e(rx/N)\right)\right|`$
$`=`$ $`{\displaystyle \frac{1}{|B|}}\left|{\displaystyle \underset{xB}{}}\left(1\mathrm{cos}(2\pi rx/N)\right)\right|`$
$``$ $`4\pi ^2\underset{xB}{sup}rx/N^2`$
$``$ $`16ϵ^2.`$
A very similar calculation shows that

$$\left|1\stackrel{~}{\beta }(2r)\right|64ϵ^2,$$

and the lemma follows quickly.

Proof of Proposition 6.4. By (6.3) we have, observing that $`\stackrel{~}{a_1}=\stackrel{~}{a}\stackrel{~}{\beta }^2`$,

$`{\displaystyle a_1(x)a_1(x+d)a_1(x+2d)}`$ $``$ $`{\displaystyle a_1(x)a_1(x+d)a_1(x+2d)}`$ (6.7)
$`{\displaystyle a(x)a(x+d)a(x+2d)}+(\mathrm{log}N)^3N^2`$
$`=`$ $`O(N^{3/2})`$
$`N^1{\displaystyle \underset{r}{}}\stackrel{~}{a}(r)^2\stackrel{~}{a}(2r)\left(1\stackrel{~}{\beta }(r)^4\stackrel{~}{\beta }(2r)^2\right).`$
Split the sum in (6.7) into two parts, that over $`rR`$ and that over $`rR`$. When $`rR`$ we use Lemma 6.7 to get

$`{\displaystyle \underset{rR}{}}\stackrel{~}{a}(r)^2\stackrel{~}{a}(2r)\left(1\stackrel{~}{\beta }(r)^4\stackrel{~}{\beta }(2r)^2\right)`$ $``$ $`2^{12}ϵ^2|R|`$
$``$ $`Cϵ^2\delta ^{5/2},`$
this last inequality following from Lemma 6.6 with $`p=5/2`$. To estimate the sum over $`rR`$, we again use Lemma 6.6 with $`p=5/2`$. Indeed using Hölder’s inequality we have

$`\left|{\displaystyle \underset{rR}{}}\stackrel{~}{a}(r)^2\stackrel{~}{a}(2r)\left(1\stackrel{~}{\beta }(r)^4\stackrel{~}{\beta }(2r)^2\right)\right|`$ $``$ $`2\underset{rR}{sup}|\stackrel{~}{a}(r)|^{1/2}{\displaystyle \underset{r}{}}|\stackrel{~}{a}(r)|^{5/2}`$
$``$ $`C\delta ^{1/2}.`$
This concludes the proof of Proposition 6.4.

By (6.6) and Lemma 6.3, $`a_1`$ behaves a bit like a measure associated to a set of size $`\alpha N`$. As promised, we use this information together with an argument originally due to Varnavides  to get a lower bound on $`a_1(x)a_1(x+d)a_1(x+2d)`$.

###### Lemma 6.8.

For some absolute constant $`C_2`$ we have

$$\underset{x,d_N}{}a_1(x)a_1(x+d)a_1(x+2d)\mathrm{exp}\left(C_2\alpha ^2\mathrm{log}(1/\alpha )\right)N^1.$$

Proof. Let $`A^{}=\{x_N:a_1(x)\alpha /2N\}`$. By Lemma 6.3 we have

$$\alpha a_1(x)\frac{2|A^{}|}{N}+\frac{\alpha }{2N}|A^c|,$$

which implies that $`|A^{}|\alpha N/4`$. We will give a lower bound for $`Z`$, the number of 3APs in $`A^{}`$. It is clear that $`a_1(x)a_1(x+d)a_1(x+2d)`$ is at least $`\alpha ^3Z/8N^3`$.

Now by Bourgain’s theorem333We could equally well use Roth’s original theorem here, at the expense of making any bounds for the relative density in Theorem 1.4 even worse.  there is a constant $`C_1`$ such that if

$$M\mathrm{exp}\left(C_1\alpha ^2\mathrm{log}(1/\alpha )\right)$$

then any subset of $`\{1,\mathrm{},M\}`$ of density at least $`\alpha /8`$ contains a 3AP with non-zero common difference. Now there are exactly $`N(N1)`$ non-trivial arithmetic progressions of length $`M`$ in $`_N`$, and $`A^{}`$ will have density at least $`\alpha /8`$ on many of them. To estimate exactly how many, fix a common difference $`d0`$, and let $`I=\{0,d,2d,\mathrm{},(M1)d\}`$. We have $`_xA^{}I(x)\alpha NM/4`$, but $`A^{}I(x)M`$ for every $`x`$. Thus another simple averaging argument shows that $`A^{}I(x)\alpha M/8`$ for at least $`\alpha N/8`$ values of $`x`$.

In total, then, there are at least $`\alpha N^2/8`$ progressions of length $`M`$ on which $`A^{}`$ has density at least $`\alpha /8`$. Each of them contains a 3AP consisting of elements of $`A^{}`$. No 3AP thus counted can arise from more than $`M^2`$ progressions of length $`M`$. Thus we have two different ways of bounding $`Z`$, and putting them together gives

$$Z\alpha N^2/8M^2.$$

The lemma follows.

Combining this with Proposition 6.4, we get

$$C^{}N^{1/2}+2^{12}ϵ^2\delta ^{5/2}+C\delta ^{1/2}\mathrm{exp}\left(C_2\alpha ^2\mathrm{log}(1/\alpha )\right).$$
(6.8)
There are constants $`C_3,C_4`$ so that if we choose

$$\delta =\mathrm{exp}\left(C_3\alpha ^2\mathrm{log}(1/\alpha )\right)$$

and

$$ϵ=\mathrm{exp}\left(C_4\alpha ^2\mathrm{log}(1/\alpha )\right)$$

then (6.8) cannot hold, and we will have derived a contradiction to the assumption that $`A`$ contains no 3APs. We are permitted to choose any values of $`ϵ`$ and $`\delta `$ so that the condition of Lemma 6.8 is satisfied. Recalling that $`kC\delta ^{5/2}`$ (a consequence of Lemma 6.6) and that $`W\mathrm{log}\mathrm{log}N/8`$, we see that (6.8) can indeed be contradicted provided that

$$\alpha C\sqrt{\frac{\mathrm{log}_5N}{\mathrm{log}_4N}}.$$
(6.9)
The subscripts indicate the number of iterated logarithms, not the base to which those logarithms are taken!

Let us remind the reader of what it is that we have contradicted. We assumed that there was a subset $`A_0𝒫`$ of positive relative upper density, containing no 3AP. The number $`\alpha `$ was related to the relative upper density of $`A_0`$, via the slightly technical reductions made in Lemma 6.1. A bound of the form (6.9) also holds for $`\alpha _0`$. That is, any subset of $`𝒫_n`$ with cardinality at least $`Cn(\mathrm{log}_5n)^{1/2}/\mathrm{log}n(\mathrm{log}_4n)^{1/2}`$ contains a 3AP.

By far the most important reason for us getting such a poor bound was the need to prove Lemma 6.2, which says that by passing to a subprogression of common difference $`m=_{pW}p`$ one can make the primes look somewhat uniform. This is a rather crude trick but we have not been able to get around it. Even if we could, the resultant bounds would surely be many miles from the probable truth, which is that any subset of $`[N]`$ of cardinality $`N(\mathrm{log}N)^{1000}`$ contains 3APs.

Let us conclude by remarking that the methods of this section use rather little about the primes. In fact by the same argument one could establish a Roth-type theorem relative to any measure $`\mu :_N^+`$ for which one had good control on $`sup_{r0}|\stackrel{~}{\mu }(r)|`$ together with bounds for $`\stackrel{~}{f}_p`$, for some $`p(2,3)`$ and any $`f`$ satisfying $`0f(x)\mu (x)`$ pointwise. In practise bounds of this latter type will come by restriction theory arguments of the type given in §5. A more general setting for our arguments, along the lines just described, is given in .

## 7. Acknowledgements

The author would like to thank Tim Gowers for his insights into Vinogradov’s three-primes theorem, which played a substantial part in the development of this paper. He would also like to thank Imre Ruzsa for helpful conversations, Jean Bourgain for drawing his attention to the references  and the students who attended the course  for their enthusiasm.


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# The Lefschetz theorem for CR submanifolds and the nonexistence of real analytic Levi flat submanifolds

## 1. Introduction

In this paper we study the topology of a CR submanifold with degenerate Levi form embedded in a compact Hermitian symmetric space. In the case that the ambient manifold is $`^v`$ we prove a Lefschetz hyperplane theorem for such CR submanifolds. In addition we relate the homotopy groups of the ambient manifold with the homotopy groups of the CR submanifold within a range depending on the nullity of the Levi form. The optimal results are obtained for Levi flat submanifolds. In particular we show that a compact Levi flat submanifold in a compact Hermitian symmetric space of complex dimension $`v`$ is simply connected if the real dimension of the Levi flat manifold is greater than $`v+1`$. We then apply this result to give a proof of the nonexistence of real analytic Levi flat submanifolds of real dimension greater than $`v+1`$. The proof of the Lefschetz theorem is based on the Morse theory of paths between a CR submanifold and a complex submanifold. The proof of the other restrictions on the topology of a CR submanifold is based on the Morse theory of paths between two CR submanifolds. The crucial ingredient in both proofs is a second variation calculation in Schoen-Wolfson \[S-W\], that goes back to Frankel \[F\].

To our knowledge results of this type relating the topology of a CR submanifold with its Levi form are the first of their kind. They are natural generalizations of analogous results for complex submanifolds. On the nonexistence of Levi flat submanifolds: Lins Neto proved the nonexistence of real analytic Levi flat hypersurfaces in $`^v`$, for $`v3`$ \[LN\]; Ohsawa \[O\] did the $`v=2`$ case, and Siu \[S\] proved the nonexistence of smooth Levi flat hypersurfaces in $`^v`$. Whether our nonexistence result remains true in the smooth category is an interesting open problem.

## 2. Definitions and the statements of the theorems

We will use the following notion defined in \[K-W\]: Let $`V`$ be an irreducible compact Kähler manifold of complex dimension $`v`$ with nonnegative holomorphic bisectional curvature. Define the symmetric bilinear form $`H_Y(W,Z)=R(Y,JY)W,JZ`$, where $`Y,W,ZTV`$. Then for any $`Y0`$, $`H_Y`$ is positive semi-definite. Denote the null-space by $`𝒩_Y`$. Let $`\mathrm{}(Y)`$ be the complementary dimension of $`𝒩_Y`$ and define the complex positivity $`\mathrm{}`$ of $`V`$ by

$$\mathrm{}=\underset{Y0}{inf}\mathrm{}(Y).$$

By the classification result of Mok \[M\], the above definition is independent of the choice of the base point. An equivalent definition is used in \[K-W\] where $`\mathrm{}`$ is computed for all compact Hermitian symmetric spaces. In particular in \[K-W\] it is shown:

$`^v`$: $`\mathrm{}=v.`$

$`\text{Gr}(p,p+q;)`$: $`\mathrm{}=p+q1.`$

$`\text{Gr}(2,p+2;)`$: $`\mathrm{}=p1.`$

$`Sp(r)/U(r)`$: $`\mathrm{}=r.`$

$`SO(2r)/U(r)`$: $`\mathrm{}=(r1)+(r2).`$

$`E_6/(Spin(10)\times T^1)`$: $`\mathrm{}=11`$

$`E_7/(E_6\times T^1)`$: $`\mathrm{}=17`$.

This computation goes back to \[B\] and \[C-V\]. It is also related to the concept called degree of nondegeneracy of the bisectional curvature in the Kähler geometry (cf. \[S\] for a definition). In fact, if $`s`$ is the degree of the nondegeneracy of $`V`$ one has $`s+\mathrm{}=v+1`$.

Let $`M`$ be a compact real submanifold of real dimension $`2p+1`$. For $`xM`$, set $`T_x^{1,0}(M)=T_x^{1,0}(V)(T_x(M))`$. Then $`T^{1,0}(M)`$ is a subbundle of $`T(M)`$. We say that $`(M,T^{1,0}(M))`$ is a CR submanifold of $`V`$ with CR structure $`T^{1,0}(M)`$ if the following are satisfied:

1. $`\mathrm{dim}_{}T^{1,0}(M)=p.`$
2. $`T^{1,0}(M)T^{0,1}(M)=\{0\}`$, where $`T^{0,1}(M)=\overline{T^{1,0}(M)}.`$
3. For any two sections $`X_1,X_2`$ of $`T^{1,0}(M)`$ defined on an open set $`UM`$, the Lie bracket $`[X_1,X_2]`$ is a section of $`T^{1,0}(M)`$ defined on $`U`$.

Let $`X_1,\mathrm{},X_p`$ be a local framing of $`T^{1,0}(M)`$ on an open set $`U`$. Then $`\overline{X_1},\mathrm{},\overline{X_p}`$ is a local framing of $`T^{0,1}(M)`$ on $`U`$. Choose a section $`T`$ of $`T(M)`$ so that $`X_1,\mathrm{},X_p,\overline{X_1},\mathrm{},\overline{X_p},T`$ gives a framing of $`T(M)`$ on $`U`$. Define a hermitian matrix $`(c_{ij})`$ by:

$$[X_i,\overline{X_j}]=c_{ij}T,\mathrm{mod}X_1,\mathrm{},X_p,\overline{X_1},\mathrm{},\overline{X_p}.$$

The hermitian form $`\tau `$ determined by $`(c_{ij})`$ is called the Levi form of the CR structure. The Levi matrix $`(c_{ij})`$ clearly depends on the choice of framing however the rank and nullity of the associated form $`\tau `$ are well-defined. Throughout this paper we will call the complex dimension of the null space of the Levi form of $`M`$ the nullity of $`M`$.

By studying the relative geometry of a pair of submanifolds of $`^v`$ consisting of a compact CR submanifold and a compact complex submanifold we prove the following Lefschetz theorems for compact CR submanifolds with degenerate Levi form.

###### Theorem 2.1.

Let $`M`$ be a compact CR submanifold of $`^v`$ of dimension $`2p+1`$. Suppose that the nullity of $`M`$ is everywhere greater than or equal to $`r`$, where $`0<rp`$. If $`jr1`$ then

$$\pi _j(M,M^{v1})=0.$$

Hence, if $`j<r1`$ then the map induced by inclusion

$$\pi _j(M^{v1})\pi _j(M),$$

is an isomorphism and if $`j=r1`$ this map is onto.

Note that, for generic $`^{v1}^v`$, $`M^{v1}`$ is a compact CR submanifold of $`^{v1}`$ of dimension $`2(p1)+1`$ with nullity everywhere greater than or equal to $`r1`$, where $`0<r1p1`$. Therefore the result can be iterated.

We say that $`M`$ is Levi flat if the Levi form of the CR structure vanishes. In particular, a Levi flat submanifold admits a codimension one foliation with complex analytic leaves.

###### Corollary 2.2.

Let $`M`$ be a compact Levi flat submanifold of $`^v`$ of dimension $`m=2p+1`$. If $`jp1`$ then

$$\pi _j(M,M^{v1})=0.$$

Hence, if $`j<p1`$ then the map induced by inclusion

$$\pi _j(M^{v1})\pi _j(M),$$

is an isomorphism and if $`j=p1`$ this map is onto.

For generic $`^{v1}`$, $`M^{v1}`$ is a compact Levi flat submanifold of $`^{v1}`$ of dimension $`2(p1)+1`$ and so, as above, the result can be iterated. Combining Corollary 2.2 with the Hurewicz isomorphism theorem we get the homology version of the Lefschetz theorem on hyperplane sections for Levi flat submanifolds of $`^v`$.

By studying the relative geometry of a pair of compact CR submanifolds of a compact hermitian symmetric space $`V`$ with dim$`{}_{}{}^{}V=v`$ we derive restrictions on the topology of a compact CR submanifold with degenerate Levi form.

###### Theorem 2.3.

Let $`V`$ be an irreducible compact Kähler manifold of complex dimension $`v`$ with nonnegative holomorphic bisectional curvature of complex positivity $`\mathrm{}`$. Let $`M`$ be a CR submanifold of real dimension $`2p+1`$ and suppose that the nullity of $`M`$ is everywhere greater than or equal to $`r`$, where $`0<rp`$. Then,

$$\pi _j(V,M)=0,$$

provided $`j\mathrm{}+12(vr)`$.

In particular if $`M`$ is Levi flat of real dimension $`m`$ and $`m2(v+1)\mathrm{}`$ then it follows that $`M`$ is simply connected. The next result then follows from a theorem of Haefliger \[H\].

###### Theorem 2.4.

Let $`V`$ be an irreducible compact Kähler manifold of complex dimension $`v`$ with nonnegative holomorphic bisectional curvature of complex positivity $`\mathrm{}`$. There are no real analytic Levi flat submanifolds of dimension $`m`$ in $`V`$ when $`m2(v+1)\mathrm{}`$.

## 3.  The relative geometry of a complex submanifold and a CR submanifold

Let $`V`$ be an irreducible compact Kähler manifold of complex dimension $`v`$, with nonnegative bisectional curvature and with complex structure $`J`$. Let $`\mathrm{}`$ be the complex positivity of $`V`$. Let $`M`$ be a real smooth compact submanifold of dimension $`2p+1`$ and let $`N`$ be a compact complex submanifold of complex dimension $`n`$.

We denote, by $`\mathrm{\Omega }(V;M,N)=\mathrm{\Omega }`$, the space of piecewise smooth paths $`\gamma :[0,1]V`$ constrained by the requirements that $`\gamma (0)M,\gamma (1)N`$. Consider the energy of a path

$$E(\gamma )=_0^1|\dot{\gamma }|^2𝑑t$$

as a function on $`\mathrm{\Omega }`$. It is shown in \[S-W\] that $`\gamma `$ is a critical point of $`E`$ if:

1. $`\gamma `$ is a smooth geodesic
2. $`\gamma `$ is normal to $`M`$ at $`\gamma (0)`$ and normal to $`N`$ at $`\gamma (1)`$, respectively.

Let $`W_1,W_2T_\gamma \mathrm{\Omega }`$. If $`\gamma `$ is a critical point of $`E`$ then the second variation of $`E`$ along $`\gamma `$ is:

$`{\displaystyle \frac{1}{2}}E_{}(W_1,W_2)`$ $`=`$ $`_{W_1}W_2,\dot{\gamma }|_0^1+{\displaystyle _0^1}_{\dot{\gamma }}W_1,_{\dot{\gamma }}W_2𝑑t`$
$`{\displaystyle _0^1}R(\dot{\gamma },W_1)\dot{\gamma },W_2𝑑t.`$
where $`R`$ denotes the curvature tensor of $`V`$.

###### Theorem 3.1.

Let $`M`$ be a CR submanifold in $`V`$ of dimension $`2p+1`$. Suppose that the nullity of $`M`$ is everywhere greater than or equal to $`r`$, where $`0<rp`$. Let $`N`$ be a compact complex submanifold of complex dimension $`n`$. Then the index of a critical point $`\gamma `$ of $`E`$ is at least $`\mathrm{}+1(vr)(vn)`$, where $`\mathrm{}`$ is the complex positivity of $`V`$.

###### Proof.

Since the null space of the Levi form of $`M`$ everywhere has dimension greater than or equal to $`r`$, for any point $`xM`$, locally there are $``$-linearly independent vector fields $`X_1,X_2,\mathrm{},X_r`$ that are sections of $`T^{1,0}(M)`$ and that lie in the null space of $`\tau `$. We can choose the $`X_i`$ so that if $`X_i=W_i\sqrt{1}JW_i`$ then $`W_1,\mathrm{},W_r,JW_1,\mathrm{},JW_r`$ are orthonormal. Now we follow the argument of \[S-W\] and estimate the index of a critical point $`\gamma `$ of $`E`$.

By the second variation formula we have that

$`{\displaystyle \frac{1}{2}}E_{}(W,W)`$ $`=`$ $`<_WW,\dot{\gamma }>|_0^1+{\displaystyle _0^1}<_{\dot{\gamma }}W,_{\dot{\gamma }}W>`$
$`{\displaystyle _0^1}<R(\dot{\gamma },W)\dot{\gamma },W>dt.`$
Parallel translate $`W_i`$, $`i=1,\mathrm{},r`$, along $`\gamma `$ to construct a vector field $`W_i(t)`$ along $`\gamma `$. Of course, $`W_i(1)`$ need not be tangent to $`N`$ at $`\gamma (1)`$ so $`W_i(t)`$ is not necessarily an element of $`T_\gamma \mathrm{\Omega }`$. However whenever $`W_i(1)TN`$ we have:

(3.1)
$$\frac{1}{2}E_{}(W_i,W_i)=<_{W_i}W_i,\dot{\gamma }>|_0^1_0^1<R(\dot{\gamma },W_i)\dot{\gamma },W_i>dt.$$

Similarly, parallel translate $`JW_i`$ along $`\gamma `$ to get:

(3.2)
$$\frac{1}{2}E_{}(JW_i,JW_i)=<_{JW_i}JW_i,\dot{\gamma }>|_0^1_0^1<R(\dot{\gamma },JW_i)\dot{\gamma },JW_i>dt.$$

Observe that,

$`<_{JW_i}JW_i,\dot{\gamma }>`$ $`=`$ $`<J(_{JW_i}W_i),\dot{\gamma }>`$
$`=`$ $`<_{JW_i}W_i,J\dot{\gamma }>`$
$`=`$ $`<_{W_i}JW_i,J\dot{\gamma }>+<[JW_i,W_i],J\dot{\gamma }>`$
$`=`$ $`<_{W_i}W_i,\dot{\gamma }>+<[JW_i,W_i],J\dot{\gamma }>,`$
At $`\gamma (0)`$ the vector $`\dot{\gamma }`$ is normal to $`M`$. Thus at $`\gamma (0)`$ the vector $`J\dot{\gamma }`$ is either normal to $`M`$ or, if tangent to $`M`$, then it is normal to $`T^{1,0}(M)T^{0,1}(M)`$. By assumption $`X_i(0)`$ lies in the null space of the Levi form, so at $`\gamma (0)`$, $`<[JW_i,W_i],J\dot{\gamma }>=0`$. To show that $`<[JW_i,W_i],J\dot{\gamma }>=0`$ at $`\gamma (1)`$ we use a simple argument of Frankel \[F\]. Let $`R`$ be an analytic curve that passes through $`\gamma (1)`$ with tangent vectors $`W_i,JW_i`$ at $`\gamma (1)`$ and extend $`W_i,JW_i`$ to be tangent vector fields to $`R`$ in a neighborhood $`U`$ of $`\gamma (1)`$. Then on $`U`$, $`[JW_i,W_i]`$ and $`J[JW_i,W_i]`$ are tangent to $`R`$ and so, at $`\gamma (1)`$, they are perpendicular to $`\dot{\gamma }(1)`$.

Therefore adding (2.1) and (2.2) we have

$`{\displaystyle \frac{1}{2}}E_{}(W_i,W_i)`$ $`+`$ $`{\displaystyle \frac{1}{2}}E_{}(JW_i,JW_i)=`$
$``$ $`{\displaystyle _0^1}(<R(\dot{\gamma },W_i)\dot{\gamma },W_i>+<R(\dot{\gamma },JW_i)\dot{\gamma },JW_i>)dt.`$
Using the symmetries of the curvature tensor we have:

(3.3)
$$R(\dot{\gamma },W_i)\dot{\gamma },W_i+R(\dot{\gamma },JW_i)\dot{\gamma },JW_i=R(\dot{\gamma },J\dot{\gamma })W_i,JW_i.$$

This expression is the holomorphic bisectional curvature of the complex lines $`\dot{\gamma }J\dot{\gamma }`$ and $`W_iJW_i`$.

Note that the vectors $`W_i(1),JW_i(1)`$ are perpendicular to both $`\dot{\gamma }(1)`$ and $`J\dot{\gamma }(1)`$. Thus the vector space

$$S=\text{span}\{W_1(1),\mathrm{},W_r(1),JW_1(1),\mathrm{},JW_r(1)\}$$

is a complex $`r`$-dimensional space lying in a complex $`(v1)`$-dimensional subspace of $`T_{\gamma (1)}V`$. It follows that the subspace $`ST_{\gamma (1)}N`$ has complex dimension at least equal to

$$r+n(v1).$$

Moreover, the vector fields $`\{W,JW\}`$ with $`W(1),JW(1)ST_{\gamma (1)}N`$ are parallel and lie in $`T_\gamma \mathrm{\Omega }`$. Therefore under the assumption that the positivity of $`V`$ is $`\mathrm{}`$ it follows that every nontrivial critical point $`\gamma `$ of $`E`$ on $`\mathrm{\Omega }`$ has index $`\lambda \lambda _0`$ with,

$$\lambda _0=r+n(v1)(v\mathrm{}).$$

Using the relation between Morse theory on the path space $`\mathrm{\Omega }`$ and the topology of $`(V;M,N)`$ (see \[S-W\] or \[K-W\] for details) we have:

###### Theorem 3.2.

Let $`V`$ be a compact Kähler manifold of non-negative holomorphic bisectional curvature, of complex dimension $`v`$ and with complex positivity $`\mathrm{}`$. Let $`M`$ be a compact CR submanifold of $`V`$ of dimension $`m=2k+1`$ and suppose that the nullity of $`M`$ is everywhere greater than or equal to $`r`$, where $`0<rk`$. Let $`N`$ be a compact complex submanifold of complex dimension $`n`$. Then the homomorphism induced by the inclusion

$$ı_{}:\pi _j(M,MN)\pi _j(V,N)$$

is an isomorphism for $`jn+rv(v\mathrm{})`$ and is a surjection for $`j=n+rv(v\mathrm{})+1`$ .

The following result is proved in \[S-W\] and \[K-W\] using techniques analogous to those used above.

###### Theorem 3.3.

Let $`V`$ be a compact Kähler manifold of non-negative holomorphic bisectional curvature, of complex dimension $`v`$ and with complex positivity $`\mathrm{}`$. Let $`N`$ be a compact complex submanifold of complex dimension $`n`$. If

$$j2nv(v\mathrm{})+1$$

then

$$\pi _j(V,N)=0.$$

###### Corollary 3.4.

Under the same hypothesis as in Theorem 3.2, if

$$j\mathrm{min}(2nv(v\mathrm{})+1,n+rv(v\mathrm{}))$$

then

$$\pi _j(M,MN)=0.$$

###### Proof.

Follows from Theorem 3.3 and Theorem 3.2. ∎

Applying Corollay 3.4 to the special case when $`V=^v`$ with the Fubini-Study metric and $`N=^{v1}`$ is a hyperplane we have the Lefschetz hyperplane theorem for compact CR submanifolds with degenerate Levi form (Theorem 2.1) and for Levi flat submanifolds (Corollary 2.2).

## 4.  The relative geometry of a pair of CR submanifolds

Let $`V`$ be an irreducible compact Kähler manifold of complex dimension $`v`$, with nonnegative bisectional curvature and with complex structure $`J`$. Let $`\mathrm{}`$ be the complex positivity of $`V`$. Let $`M`$ and $`N`$ be a real smooth compact submanifolds of dimensions $`m=2p+1`$ and $`n=2q+1`$, respectively.

We denote, by $`\mathrm{\Omega }(V;M,N)=\mathrm{\Omega }`$, the space of piecewise smooth paths $`\gamma :[0,1]V`$ constrained by the requirements that $`\gamma (0)M,\gamma (1)N`$. Consider the energy of a path

$$E(\gamma )=_0^1|\dot{\gamma }|^2𝑑t$$

as a function on $`\mathrm{\Omega }`$. It is shown in \[S-W\] that $`\gamma `$ is a critical point of $`E`$ if:

1. $`\gamma `$ is a smooth geodesic
2. $`\gamma `$ is normal to $`M`$ at $`\gamma (0)`$ and normal to $`N`$ at $`\gamma (1)`$, respectively.

Let $`W_1,W_2T_\gamma \mathrm{\Omega }`$. If $`\gamma `$ is a critical point of $`E`$ then the second variation of $`E`$ along $`\gamma `$ is:

$`{\displaystyle \frac{1}{2}}E_{}(W_1,W_2)`$ $`=`$ $`_{W_1}W_2,\dot{\gamma }|_0^1+{\displaystyle _0^1}_{\dot{\gamma }}W_1,_{\dot{\gamma }}W_2𝑑t`$
$`{\displaystyle _0^1}R(\dot{\gamma },W_1)\dot{\gamma },W_2𝑑t.`$
where $`R`$ denotes the curvature tensor of $`V`$.

###### Theorem 4.1.

Let $`M`$ and $`N`$ be CR submanifolds of $`V`$ of dimensions $`2p+1`$ and $`2q+1`$, respectively. Suppose that the nullity of $`M`$ and $`N`$ are everywhere greater than or equal to $`r`$ and $`s`$, respectively, where $`0<rp`$ and $`0<sq`$. Then the index of a critical point $`\gamma `$ of $`E`$ is at least $`\mathrm{}+1(vr)(vs)`$, where $`\mathrm{}`$ is the complex positivity of $`V`$.

###### Proof.

Since the null space of the Levi form of $`M`$ everywhere has dimension greater than or equal to $`r`$, for any point $`xM`$, locally there are $``$-linearly independent vector fields $`X_1,X_2,\mathrm{},X_r`$ that are sections of $`T^{1,0}(M)`$ and that lie in the null space of the Levi form of $`M`$. We can choose the $`X_i`$ so that if $`X_i=W_i\sqrt{1}JW_i`$ then $`W_1,\mathrm{},W_r,JW_1,\mathrm{},JW_r`$ are orthonormal. Now we follow the argument of \[S-W\] and estimate the index of a critical point $`\gamma `$ of $`E`$.

By the second variation formula we have that

$`{\displaystyle \frac{1}{2}}E_{}(W,W)`$ $`=`$ $`<_WW,\dot{\gamma }>|_0^1+{\displaystyle _0^1}<_{\dot{\gamma }}W,_{\dot{\gamma }}W>`$
$`{\displaystyle _0^1}<R(\dot{\gamma },W)\dot{\gamma },W>dt.`$
Parallel translate $`W_i`$, $`i=1,\mathrm{},r`$, along $`\gamma `$ to construct a vector field $`W_i(t)`$ along $`\gamma `$. Of course, $`W_i(1)`$ need not be tangent to $`N`$ at $`\gamma (1)`$ so $`W_i(t)`$ is not necessarily an element of $`T_\gamma \mathrm{\Omega }`$. However whenever $`W_i(1)TN`$ we have:

(4.1)
$$\frac{1}{2}E_{}(W_i,W_i)=<_{W_i}W_i,\dot{\gamma }>|_0^1_0^1<R(\dot{\gamma },W_i)\dot{\gamma },W_i>dt.$$

Similarly, parallel translate $`JW_i`$ along $`\gamma `$ to get:

(4.2)
$$\frac{1}{2}E_{}(JW_i,JW_i)=<_{JW_i}JW_i,\dot{\gamma }>|_0^1_0^1<R(\dot{\gamma },JW_i)\dot{\gamma },JW_i>dt.$$

Now assume that $`X_i(1)=W_i(1)\sqrt{1}JW_i(1)`$ lies in the null space of the Levi form of $`N`$ at $`\gamma (1)`$. Observe that,

$`<_{JW_i}JW_i,\dot{\gamma }>`$ $`=`$ $`<J(_{JW_i}W_i),\dot{\gamma }>`$
$`=`$ $`<_{JW_i}W_i,J\dot{\gamma }>`$
$`=`$ $`<_{W_i}JW_i,J\dot{\gamma }>+<[JW_i,W_i],J\dot{\gamma }>`$
$`=`$ $`<_{W_i}W_i,\dot{\gamma }>+<[JW_i,W_i],J\dot{\gamma }>,`$
At $`\gamma (0)`$ the vector $`\dot{\gamma }`$ is normal to $`M`$. Thus at $`\gamma (0)`$ the vector $`J\dot{\gamma }`$ is either normal to $`M`$ or, if tangent to $`M`$, then it is normal to $`T^{1,0}(M)T^{0,1}(M)`$. By assumption $`X_i(0)`$ lies in the null space of the Levi form, so at $`\gamma (0)`$, $`<[JW_i,W_i],J\dot{\gamma }>=0`$. To show that $`<[JW_i,W_i],J\dot{\gamma }>=0`$ at $`\gamma (1)`$ we use the assumption that $`X_i(1)`$ lies in the null space of the Levi form of $`N`$ at $`\gamma (1)`$ and apply the previous reasoning. Therefore adding (2.1) and (2.2) we have

$`{\displaystyle \frac{1}{2}}E_{}(W_i,W_i)`$ $`+`$ $`{\displaystyle \frac{1}{2}}E_{}(JW_i,JW_i)=`$
$``$ $`{\displaystyle _0^1}(<R(\dot{\gamma },W_i)\dot{\gamma },W_i>+<R(\dot{\gamma },JW_i)\dot{\gamma },JW_i>)dt.`$
Using the symmetries of the curvature tensor we have:

(4.3)
$$R(\dot{\gamma },W_i)\dot{\gamma },W_i+R(\dot{\gamma },JW_i)\dot{\gamma },JW_i=R(\dot{\gamma },J\dot{\gamma })W_i,JW_i.$$

This expression is the holomorphic bisectional curvature of the complex lines $`\dot{\gamma }J\dot{\gamma }`$ and $`W_iJW_i`$.

Note that the vectors $`W_i(1),JW_i(1)`$ are perpendicular to both $`\dot{\gamma }(1)`$ and $`J\dot{\gamma }(1)`$. Thus the vector space

$$S=\text{span}\{W_1(1),\mathrm{},W_r(1),JW_1(1),\mathrm{},JW_r(1)\}$$

is a complex $`r`$-dimensional space lying in a complex $`(v1)`$-dimensional subspace of $`T_{\gamma (1)}V`$. Denote the null space of the Levi form of $`N`$ at $`xN`$ by $`𝒵_x(TN)`$. Then the subspace $`S𝒵_{\gamma (1)}(TN)`$ has complex dimension at least equal to

$$r+s(v1).$$

Moreover, the vector fields $`\{W,JW\}`$ with $`W(1),JW(1)S𝒵_{\gamma (1)}(TN)`$ are parallel and lie in $`T_\gamma \mathrm{\Omega }`$. Therefore under the assumption that the positivity of $`V`$ is $`\mathrm{}`$ it follows that every nontrivial critical point $`\gamma `$ of $`E`$ on $`\mathrm{\Omega }`$ has index $`\lambda \lambda _0`$ with,

$$\lambda _0=r+s(v1)(v\mathrm{}).$$

Using the relation between Morse theory on the path space $`\mathrm{\Omega }`$ and the topology of $`(V;M,N)`$ (see \[S-W\] or \[K-W\] for details) we have:

###### Theorem 4.2.

Let $`V`$ be a compact Kähler manifold of non-negative holomorphic bisectional curvature, of complex dimension $`v`$ and with complex positivity $`\mathrm{}`$. Let $`M,N`$ be a compact CR submanifolds of $`V`$ of dimensions $`2p+1`$ and $`2q+1`$, respectively and suppose that the nullity of $`M`$ and $`N`$ are everywhere greater than or equal to $`r`$ and $`s`$, respectively, where $`0<rp`$ and $`0<sq`$. Then the homomorphism induced by the inclusion

$$ı_{}:\pi _j(M,MN)\pi _j(V,N)$$

is an isomorphism for $`jr+sv(v\mathrm{})`$ and is a surjection for $`j=r+sv(v\mathrm{})+1`$ .

Applying this result to the case $`M=N`$ we have:

###### Corollary 4.3.

Let $`M`$ be a CR submanifold of $`V`$ as in Theorem 4.2 . Then

$$\pi _j(V,M)=0,$$

provided $`j\mathrm{}+12(vr)`$, where $`\mathrm{}`$ is the complex positivity of $`V`$.

Using the exact homotopy sequence of the pair $`(V,M)`$ and that $`\pi _1(V)=0`$ we conclude:

###### Corollary 4.4.

Let $`V`$ and $`M`$ be as above. If $`2r2v+1\mathrm{}`$ then $`\pi _1(M)=\{0\}`$.

###### Corollary 4.5.

Let $`V`$ be as above and suppose that $`M`$ is Levi flat. If dim $`M2(v+1)\mathrm{}`$ then $`\pi _1(M)=\{0\}`$.

Recall the well known result of A. Haefliger \[H\].

###### Theorem 4.6.

There is no real analytic, codimension one foliation of a compact simply-connected manifold.

Corollary 4.5 and Theorem 4.6 imply Theorem 2.4.

Combining the topological restrictions imposed on Levi flat submanifolds by Theorem 4.2 and its corollaries with codimension one foliation theory leads to many partial nonexistence results such as:

###### Theorem 4.7.

Let $`V`$ be an irreducible compact Kähler manifold of complex dimension $`v`$ with complex positivity $`\mathrm{}`$. If $`m2(v+1)\mathrm{}`$ then there are no smooth Levi flat submanifolds of dimension $`m`$ in $`V`$ with a compact leaf.

We leave the formulation of similar results to the reader.


text:
# Solution of scaling quantum networks

## Abstract

We show that all scaling quantum graphs are explicitly integrable, i.e. any one of their spectral eigenvalues $`E_n`$ is computable analytically, explicitly, and individually for any given $`n`$. This is surprising, since quantum graphs are excellent models of quantum chaos \[see, e.g., T. Kottos and H. Schanz, Physica E 9, 523 (2001)\].

Graphs are networks of bonds and vertices. Figure 1 shows two examples: a three-bond four-vertex star graph (Fig. 1a) and a three-bond four-vertex linear graph (Fig. 1b). A quantum particle moving on the graph turns the graph into a quantum graph . If the quantum particle moves freely on the graph, subjected only to flux conservation at its vertices, we call it a standard quantum graph. This is the type of quantum graph most frequently studied in the literature . A larger class of quantum graphs, including the standard quantum graphs, are dressed quantum graphs . A dressed quantum graph has potentials on its bonds and $`\delta `$-functions on its vertices. The potentials on its bonds are essentially arbitrary as long as they do not introduce turning points on the bonds. But even this case can be dealt with trivially by re-defining the topology of the graph.

An important subset of dressed quantum graphs are scaling quantum graphs . In this case the graph bonds are dressed with scaling potentials and the graph vertices are dressed with scaling $`\delta `$-functions. A scaling potential is one whose strength $`V_0`$ scales with the energy $`E`$ of the quantum particle according to $`V_0=\lambda E`$, where $`\lambda `$ is a constant. The strength of a scaling $`\delta `$-function scales with $`k=\sqrt{E}`$. Scaling potentials and $`\delta `$-functions are a natural choice to consider. On the one hand they frequently occur in physical systems , on the other hand they are mathematically convenient, since they allow studying a quantum system without causing phase-space metamorphoses  in the underlying classical system. It has been pointed out before  that this is the most natural way of studying quantum systems, in particular quantum chaos . Since quantum graphs are popular and successful models of quantum chaos , it may come as a surprise that the energy spectrum $`E_n`$, $`n=1,2,\mathrm{}`$ of all scaling quantum graphs is explicitly and analytically solvable in the form $`E_n=\mathrm{}`$, involving only known quantities on the right-hand side. In many cases the solutions can be stated in closed analytical form.

The spectral function $`g^{(0)}(k)`$ of a general scaling, dressed quantum graph is of the form

$$g^{(0)}(k)=\mathrm{cos}(S_0k\pi \gamma _0)\underset{j=1}{\overset{N}{}}a_j^{(0)}\mathrm{cos}(S_jk\pi \gamma _j),$$
(1)
where $`S_0>0`$ is the total reduced action length of the graph , $`0<S_j<S_0`$ are certain combinations of the reduced bond actions, $`N`$ is the number of action combinations in (1), $`\gamma _0`$, $`\gamma _j`$ are constant phases and $`a_j^{(0)}`$ are constant amplitudes. The spectrum $`E_n`$ of the quantum graph is obtained by solving the spectral equation

$$g^{(0)}(k_n^{(0)})=0,n=1,2,\mathrm{}$$
(2)
via $`E_n=(k_n^{(0)})^2`$. For the purposes of this paper we are only interested in the positive solutions of (2), and obtain a well-defined counting index $`n`$ by defining $`k_1^{(0)}`$ to be the first positive root of (2). As a first step toward the solution of the general problem, it was shown in  that (2) can be solved explicitly in the form $`k_n^{(0)}=\mathrm{}`$ if the regularity condition

$$\underset{j=1}{\overset{N}{}}|a_j^{(0)}|<1$$
(3)
is fulfilled. In order to substantiate our claim that (2) is solvable explicitly for all scaling quantum graphs, we have to show that (2) is solvable explicitly even if (3) is not fulfilled.

Before we turn our attention to the general case, we introduce our methods with the help of a simple example. Let us consider a scaling quantum graph derived from the three-bond star graph shown in Fig. 1a by putting the scaling potentials $`V_l(E)=\lambda _lE`$, $`0<\lambda _l<1`$ on its three bonds of length $`L_l`$, $`l=1,2,3`$, require the “Kirchhoff-type”  flux conservation condition $`_{l=1}^3d\psi _l/dx_l=0`$ at its central vertex ($`\psi _l`$ is the quantum wave function on bond number $`l`$ of the graph and $`x_l`$ is the coordinate on bond number $`l`$) and require Dirichlet boundary conditions on its three dead-end vertices. The spectral equation is of the form (1) with $`N=3`$, $`\gamma _1=\gamma _2=\gamma _3=0`$ and

$$S_0=\alpha _1+\alpha _2+\alpha _3,S_1=\alpha _1+\alpha _2+\alpha _3,S_2=\alpha _1\alpha _2+\alpha _3,S_3=\alpha _1+\alpha _2\alpha _3,$$
(4)
$$a_1^{(0)}=\frac{\beta _1+\beta _2+\beta _3}{\beta _1+\beta _2+\beta _3},a_2^{(0)}=\frac{\beta _1\beta _2+\beta _3}{\beta _1+\beta _2+\beta _3},a_3^{(0)}=\frac{\beta _1+\beta _2\beta _3}{\beta _1+\beta _2+\beta _3},$$
(5)
where

$$\alpha _l=\beta _lL_l,\beta _l=\sqrt{1\lambda _l},l=1,2,3.$$
(6)
The amplitudes in (5) do not fulfill the regularity condition (3). In some cases $`_{j=1}^3|a_j^{(0)}|=1`$ (for instance for $`a_j^{(0)}>0`$, $`j=1,2,3`$), and in many cases $`_{j=1}^3|a_j^{(0)}|>1`$, which strongly violates the regularity condition (3). Since the methods and techniques presented in  for obtaining the spectrum of a graph explicitly depend crucially on (3), it seems that completely different methods have to be developed for general graphs, such as the three-bond star graph of Fig. 1a, which do not fulfill (3). There is, however, a way to reduce (1) to a form that allows to bring the powerful theory of regular quantum graphs  to bear. In order to motivate and to illustrate this method, let us study the case $`\alpha _1=1`$, $`\alpha _2=7`$, $`\alpha _3=11`$, $`\beta _1=1/10`$, $`\beta _2=1/5`$, $`\beta _3=1/2`$. In this case $`a_1^{(0)}=3/4`$, $`a_2^{(0)}=1/2`$, $`a_3^{(0)}=1/4`$, $`S_0=19`$, $`S_1=17`$, $`S_2=5`$, $`S_3=3`$ and the spectral equation is given by

$$g^{(0)}(k)=\mathrm{cos}(19k)\frac{3}{4}\mathrm{cos}(17k)\frac{1}{2}\mathrm{cos}(5k)+\frac{1}{4}\mathrm{cos}(3k).$$
(7)
Since $`|a_1^{(0)}|+|a_2^{(0)}|+|a_3^{(0)}|=3/2>1`$, this quantum graph is certainly not regular. But let us look at the first derivative of (7). Dividing by $`S_0`$, we obtain

$$g^{(1)}(k)=\mathrm{cos}[S_0k+\pi /2]\underset{j=1}{\overset{3}{}}a_j^{(1)}\mathrm{cos}[S_jk+\pi /2]=$$

$$\mathrm{sin}(19k)+\frac{51}{76}\mathrm{sin}(17k)+\frac{5}{38}\mathrm{sin}(5k)\frac{3}{76}\mathrm{sin}(3k).$$
(8)
This time we have $`_{j=1}^3|a_j^{(1)}|=16/19<1`$ and therefore, since (8) is precisely of the form (1) and satisfies (3), it can be solved explicitly using the methods of . In particular it was shown in  that root number $`n`$ of a spectral equation that satisfies (3), such as (8), is found in the root interval $`[\widehat{k}_{n1},\widehat{k}_n]`$, where $`\widehat{k}_n`$ are the root separators . It was also shown in  that the location of the root separators is entirely controlled by the local extrema of the trigonometric function with the largest action argument. Thus, in our case, the root separators of (8) are given by $`\widehat{k}_n=(2n+1)\pi /38`$. Since according to  root number $`n`$ and only root number $`n`$ is located in the interval $`[\widehat{k}_{n1},\widehat{k}_n]`$, we can now compute all roots of (8) explicitly and individually according to

$$k_n^{(1)}=_{\widehat{k}_{n1}}^{\widehat{k}_n}k\left|\frac{dg^{(1)}(k)}{dk}\right|\delta (g^{(1)}(k))𝑑k.$$
(9)
In  we show that because of the hermiticity of the spectral eigenvalue problem on quantum graphs the locations of the local extrema of $`g^{(0)}(k)`$ are separators for the roots of $`g^{(0)}(k)`$. The location of the local extrema of $`g^{(0)}(k)`$, however, are given by the zeros of $`g^{(1)}(k)`$, which, up to constants, is the derivative of $`g^{(0)}(k)`$. Therefore, using the roots $`k_n^{(1)}`$, explicitly computed in (9), as the root separators of (7), we obtain, again explicitly and individually,

$$k_n^{(0)}=_{k_{n1}^{(1)}}^{k_n^{(1)}}k\left|\frac{dg^{(0)}(k)}{dk}\right|\delta (g^{(0)}(k))𝑑k.$$
(10)
This solves the task of computing the spectrum of our example of the three-bond dressed star graph explicitly.

In general, given a spectral equation (1) which does not fulfill (3), we generate a chain of derivative spectral equations $`g^{(m)}(k)`$, where $`g^{(m)}(k)`$ is the $`m`$’th derivative of (1) divided by $`S_0^m`$, explicitly given by

$$g^{(m)}(k)=\mathrm{cos}(S_0k\pi \gamma _0+m\pi /2)\underset{j=1}{\overset{N}{}}a_j^{(m)}\mathrm{cos}(S_jk\pi \gamma _j+m\pi /2),$$
(11)
where $`a_j^{(m)}=a_j^{(0)}(S_j/S_0)^m`$. Since $`S_0<S_j`$, there always exists an $`M`$ such that the amplitudes $`a_j^{(m)}`$ satisfy the regularity condition (3), i.e., $`_{j=1}^N|a_j^{(M)}|<1`$. Therefore, according to , root separators $`\widehat{k}_n^{(M)}`$ exist on the level $`M`$ and the roots $`k_n^{(M)}`$ of $`g^{(M)}(k)=0`$ are explicitly computable via

$$k_n^{(M)}=_{\widehat{k}_{n1}^{(M)}}^{\widehat{k}_n^{(M)}}k\left|\frac{dg^{(M)}(k)}{dk}\right|\delta (g^{(M)}(k))𝑑k.$$
(12)
Since we now know the roots on the level $`M`$, we can go one step backwards to level $`M1`$. According to a root-counting argument  based on the Weyl formula , the root separators $`\widehat{k}_n^{(M1)}`$ on the level $`M1`$ are the locations of the local extrema of $`g^{(M1)}(k)`$, which are given explicitly by the roots $`k_n^{(M)}`$, which we know. Therefore, $`\widehat{k}_n^{(M1)}=k_n^{(M)}`$ and the roots of $`g^{(M1)}(k)=0`$ can now be computed explicitly, according to

$$k_n^{(M1)}=_{k_{n1}^{(M)}}^{k_n^{(M)}}kS_0\left|g^{(M)}(k)\right|\delta (g^{(M1)}(k))𝑑k.$$
(13)
Steps (12) and (13) define a recursive procedure,

$$k_n^{(m1)}=_{k_{n1}^{(m)}}^{k_n^{(m)}}kS_0\left|g^{(m)}(k)\right|\delta (g^{(m1)}(k))𝑑k,m=M,M1,\mathrm{},2,$$
(14)
which can be followed until the level 0 is reached and the roots $`k_n^{(0)}`$, i.e. the spectrum of the quantum graph, is known explicitly.

It is important to notice that (12) – (14) are not just formal solutions. They yield $`k_n^{(m)}`$, $`m=0,\mathrm{},M`$ explicitly, by quadratures. Thus (12) – (14) constitute explicit solutions of the problem, very much in the spirit of the definition of explicit solutions by quadratures in the theory of differential equations .

Several special cases require discussion. If the regularity condition (3) is fulfilled, a root $`k_n`$ lies strictly inside of the interval $`[\widehat{k}_{n1},\widehat{k}_n]`$. However, if (3) is not fulfilled, it is possible that a root $`k_n^{(m)}`$ coincides with one of its separators $`\widehat{k}_{n1}^{(m)}`$ or $`\widehat{k}_n^{(m)}`$. This is, e.g., the case in our star-graph example above, where $`k_{18}^{(0)}=\pi `$ is a root and a root separator of (7). In the parameter space of $`\alpha `$’s and $`\beta `$’s cases like this are extremely rare (nongeneric). But even if such a case occurs, it does not present a problem for our theory. At the contrary, it saves one integration step since it can always be checked before performing the integration in (14), whether one of the separators $`k_{n1}^{(m)}`$ or $`k_n^{(m)}`$ is a root of $`g^{(m1)}(k)`$. If so, the result $`k_{n1}^{(m1)}=k_{n1}^{(m)}`$, or $`k_{n1}^{(m1)}=k_n^{(m)}`$, respectively, is known in advance, without actually performing the integration.

In other special cases the roots of $`g^{(m)}(k)=0`$ can be obtained in the form of explicit periodic orbit expansions . In order to illustrate this, let us return to our example of the three-bond star graph. We notice that the spectral equation $`g^{(1)}(k)=0`$ of the three-bond star graph looks the same as the spectral equation

$$g_{4\mathrm{V}\mathrm{chain}}^{(0)}(k)=\mathrm{sin}(S_0k)+r_2\mathrm{sin}(S_1k)+r_2r_3\mathrm{sin}(S_2k)r_3\mathrm{sin}(S_3k)=0$$
(15)
of the dressed four-vertex chain graph shown in Fig. 1b, where $`r_2=(\beta _1\beta _2)/(\beta _1+\beta _2)`$, $`r_3=(\beta _2\beta _3)/(\beta _2+\beta _3)`$ are the reflection coefficients at the vertices number 2 and 3 of the chain graph, and the actions $`S_0,\mathrm{},S_3`$ are the same as in (4). If we arrange for the bond actions of the chain graph to equal the bond actions of the three-bond star graph, and furthermore arrange for $`a_1^{(1)}=r_2`$, $`a_2^{(1)}=r_2r_3`$, $`a_3^{(1)}=r_3`$, which is possible if the scaling constants of the three-bond star graph fulfill $`\beta _1^2\beta _2^2+\beta _3^2=0`$, then $`g^{(1)}(k)`$ of the three-bond star graph is the same as the spectral equation (15) of the associated four-vertex chain graph and the spectral points $`k_n^{(1)}`$ can be stated immediately and explicitly in the form of convergent, periodic orbit expansions , bypassing any integrations that would have been necessary according to the scheme defined in (14).

Although they presented the first examples of explicitly solvable quantum graphs, a major shortcoming of  and  is the fact that the theory presented in  and  is only applicable to regular quantum graphs, i.e. quantum graphs that fulfill the regularity condition (3). In this paper we showed that the restriction to regular quantum graphs is not necessary: all scaling quantum graphs can be solved explicitly. Nevertheless, the theory presented in  provides an indispensable foundation without which the present theory would not be possible.

A conceptual advance is the following. Frequently an operational definition of quantum chaos, or a quantum chaotic regime, is the “loss of quantum numbers”. To illustrate, let us consider a Hamiltonian system with Hamiltonian $`\widehat{H}=\widehat{H}_0+\mu \widehat{V}`$, where $`\widehat{H}_0`$ is an integrable Hamiltonian, $`\mu `$ is a real parameter and $`\widehat{V}`$, with respect to and in conjunction with $`\widehat{H}_0`$, is a nonintegrable perturbation. Many quantum systems, for instance the hydrogen atom in a strong magnetic field , can be described in this way. For $`\mu =0`$ the system is integrable and possesses a complete set of quantum numbers

that can be obtained, at least approximately, using EBK quantization . As the parameter $`\mu `$ increases, EBK quantization breaks down and the system makes a transition to quantum chaos. This explains the frequently employed practice of characterizing the onset of quantum chaos by a loss of quantum numbers, since the breakdown of the EBK quantization scheme implies the loss of quantum numbers. The results obtained in this paper, however, show that this is not necessarily a good way to characterize quantum chaos. Although not strictly chaotic in the classical limit (due to ray-splitting  the term stochastic may characterize the situation better), quantum graphs were shown by many authors  to be excellent models of quantum chaos. Yet, our results above show that a well-defined quantum number, the counting index $`n`$, still exists, and produces explicit energy levels in exactly the same spirit as the EBK quantization scheme.

The iteration scheme (14) is perhaps the most interesting feature of our method of explicitly solving quantum graphs. We call the smallest $`M`$ that “regularizes” a given quantum graph (i.e. the amplitudes of $`g^{(M)}`$ fulfill (3)), the order of the quantum graph. For any given quantum graph its order is unique. Since the order $`M`$ of a quantum graph determines the length of the bootstrapping iteration scheme (14), it is possible that the order of a quantum graph is also an indication of the complexity of its spectrum. Quantum iterations similar to (14) were studied before  and were found to lead to sensitivity and chaos on the quantum level. This may explain the reason why certain quantum graphs are such good models of quantum chaos  and the order $`M`$ of the quantum graph may be an indication of how well a given quantum graph can be described in terms of the usual diagnostic tools of quantum chaos, such as, e.g., random matrix theory .

The authors acknowledge financial support by the National Science Foundation under Grant No. 9984075.

Figure Captions

Fig. 1: (a) Dressed three-bond star graph and (b) dressed four-vertex chain graph. Different potential strengths on the bonds are indicated by different thickness of the bonds. Different vertex strengths are indicated by different dot sizes representing the vertices.