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# Selection Criterion for Log-Linear Models Using Statistical Learning Theory

## 1 INTRODUCTION

Suppose a scientist is interested in the relation between $`n`$ features described by the random variables $`X_1,X_2,\mathrm{},X_n`$. Observing dependencies among these features is an important part of all scientific disciplines. The dependencies may be deterministic, e.g.,

$$X_n=f(X_1,\mathrm{},X_{n1})$$

or (in the generic case) statistical, i.e., the probability for the event $`X_1=x_1,X_2=x_2,\mathrm{},X_n=x_n`$, for short $`𝐗=𝐱`$, is not the product of the probabilities of the $`n`$ events $`X_j=x_j`$. All information about these dependencies is contained in the joint probability distribution $`P`$ that assigns the probability $`P(𝐱)=P(x_1,x_2,\mathrm{},x_n)`$ to each $`n`$-tuple. Notice that for large $`n`$ no “reasonable” sample size is sufficient to determine the probabilities of all $`P(x_1,x_2,\mathrm{},x_n)`$ with good reliability since the size of $`\mathrm{\Omega }:=\times _j\mathrm{\Omega }_j`$ (where each $`\mathrm{\Omega }_j`$ denotes the set of possible values $`x_j`$) is exponentially large in $`n`$. This shows that we should not try to learn the joint probability distribution, we rather have to develop inference rules for learning the statistical dependencies in a weaker sense. The task is therefore to estimate some properties of the joint distribution described by empirical data of a given size.

In the last years methods to judge scientific theories by their predicted models on real world data sets become more and more important (Hagenaars 1994, Ishii-Kuntz 1997, Pitt and Myung 2002) and since the accessible computational power is increasing strongly this trend will proceed in future.

It is a well-known fact, see e.g. (Akaike 1973, Hansen and Yu 2001), that criteria describing the goodness of fit for the model to the data set only are not enough to judge the scientific relevance of the model. Additionally one has to take the complexity of the model into account. Usually the number of free parameters are treated as a measure for the model complexity and are incooperated in criteria like in Akaike’s Information Criterion (AIC) (Akaike 1973), or Bayesian Information Criterion of Schwarz (1978). Bayesian Model Selection (Kass and Raftery 1995) and Minimum Description Length (Rissanen 1996, Hansen and Yu 2001) take also the functional form of the model into account. However, these criteria do not give any statements on the reliability of the estimation.

In this article we derive reliability statements from statistical learning theory (Vapnik 1998) and show that the trade-off between goodness of fit and model complexity can be treated with structural risk minimization.

Let us describe roughly the link between the ideas presented here and the question of minimal description length in information theory, see (Barron 1991, Rissanen 1996) and references therein. In the latter approach a risk functional is minimized over a set of probability distributions where each probability distribution is penalized according to its description length. More explicitly, the risk functional consists of two terms: the “empirical risk” evaluates the “goodness of fit” and a regularization term penalizes the code length. Barron and Cover (1991) obtain bounds on the convergence rate of the risk functional in the Hellinger distance. The difference to our approach lies in the fact the regularization function can in principle be chosen arbitrarily, only Kraft’s inequality most hold; whereas we use a penalty term coming directly from the complexity of the considered class of log-linear models and is based on statistical learning theory. Vapnik and Chervonenkis (1971) addressed the question of model selection from empirical data in a probabilistic setting (see also Vapnik (1999) and references therein). They introduced the so–called VC dimension of a function set and showed that the VC dimension of the function set which the model can implement is a crucial quantity to describe its complexity. In general the VC dimension does not agree with the usual dimension or the number of free parameters.

From the statistical learning point of view, a scientific theory delivers a prior on the class of models under consideration. Then the model is chosen according to some criterion measuring the fit on the empirical data. In this work we consider the class of all log-linear models and prior models with few statistical dependencies among the random variables, i.e. we prefer less complex models in terms of statistical independency as long as they describe the data well enough. With this assumption we estimate the model (complexity) directly from the empirical data rather than testing one model against another using test methods like the Pearson or likelihood ratio chi-square criterion, see e.g. Christensen 1997, Goodman 1978. Akaike (1973) already formulated this view point, however he did not have the conceptual tools of learning theory (Vapnik 1995).

The article is organized as follows. In Section 2 we introduce log-linear models and explain the prior for model selection leading to a risk functional on the data. We discuss the commonly used model selection criteria and explain their shortcomings. We present the idea of Markov networks and relate them to log-linear models. In Section 3.1 we calculate the VC dimension of the models under consideration and give the central result about the estimation of the risk functional for the true probability distribution. In Section 3.2 we use the idea of the structural risk minimization to formulate the model estimation for log-linear models. In Section 3.3 we show that the optimal model is the solution of a convex optimization problem.

## 2 LOG-LINEAR MODELS

### 2.1 The Empirical Risk Functional

In the work presented here we restrict our attention to discrete random variables and assume that $`𝐗`$ takes values in a finite set $`\mathrm{\Omega }=\times _{j=1}^n\mathrm{\Omega }_j`$. Assume a scientist tells us that he has found that the true joint probability distribution $`P`$ is given by $`P_t`$ and we would like to test whether he is right or not. Suppose that his model is given by the distribution that all $`n`$-tuples in $`\mathrm{\Omega }`$ occur with equal probability. Assume, in contrast, that the true distribution $`P`$ assigns the probability $`1/k`$ to $`k`$ specific $`n`$-tuples $`𝐱_1,𝐱_2,\mathrm{},𝐱_k\mathrm{\Omega }`$ and that this $`k`$ $`n`$-tuples are chosen without any simple law. Assume furthermore that $`k`$ is large compared to the available sample size but small compared to the size of $`\mathrm{\Omega }`$. Then we have only little chance to recognize that his model is wrong. Only if the $`k`$ $`n`$-tuples are selected by a simple law we would mistrust the scientist’s hypothesis. This example illustrates that we need a quality criterion for models which can be tested on a reasonable sample size. We suggest to use following criterion. For each $`n`$ tuple $`𝐱`$ in our data set we give penalty points depending on how likely its appearance was according to the hypothetical joint distribution. A good choice, for instance, is the negative logarithm of the hypothetical probability $`P_t(𝐱)`$. Let $`T`$ be the training set, then we obtain

$$\frac{1}{\mathrm{\#}T}\underset{𝐱T}{}\mathrm{ln}P_t(𝐱)$$

as the “penalty” for the model $`P_t(𝐱)`$. Notice that this sum converges in the large sample size limit to

$$R(P_t):=\underset{𝐱\mathrm{\Omega }}{}P(𝐱)\mathrm{ln}P_t(𝐱).$$
(1)
Note that $`R(P_t)`$ is closely related to the Kullback-Leibler distance (Cover and Thomas 1991) and can never be smaller than the entropy

$$S(P):=\underset{𝐱\mathrm{\Omega }}{}P(𝐱)\mathrm{ln}P(𝐱).$$

Hence, even if the scientist has predicted the true measure exactly, he will never get less than $`S(P)`$ penalty points. If $`S(P)`$ is high, he is not punished for building wrong models but for selecting the wrong features, i.e. features with too few dependencies. If the scientist considers additional random variables (“features influencing the considered ones”) more specific statements might be possible. Reasonable research consists not only in observing dependencies among a given set of features, it consists also in selecting statistically relevant features yielding a low entropy of the joint distribution. Therefore the “empirical risk functional”

$$R_{emp}(P_t):=\frac{1}{\mathrm{\#}T}\underset{𝐱T}{}\mathrm{ln}P_t(𝐱)$$
(2)
is a good measure to test the quality of the distribution $`P_t`$ for large sample size. This risk functional is well-known for estimating probabilities (Vapnik 1998).

If the space of all probability distributions is too large compared to the sample size, we should not minimize the empirical risk over all joint distributions. In order to infer a probability distribution from the statistical data we give the space a class structure. We use log-linear models (Christensen 1997, Goodman 1978) to define suitable classes. The simplest class of log-linear models is given in the case when all random variables are statistically independent. Then the logarithm of the joint probability distribution can be written as a sum of functions, each function depends on one variable only

$$\mathrm{ln}P_t(x_1,x_2,\mathrm{},x_n)=\underset{j=n}{\overset{n}{}}f_j(x_j),$$

with $`f_j(x_j)=\mathrm{ln}P(x_j)`$. The next higher class is given by allowing two-variable interaction terms

$$\mathrm{ln}P_t(x_1,\mathrm{},x_n)=\underset{jn}{}f_j(x_j)+\underset{i,jn}{}g_{i,j}(x_i,x_j).$$

The next higher class contains terms with three-variable interactions, etc.. Probability distributions with the property that their logarithm can be expressed as sum of functions with $`k`$ variables are called $`k`$-factor distributions. In other words the data is described by a $`k`$–factor log-linear model. The idea of log-linear models come from the fact that some variables influence each other directly, expressed by the interaction terms, whereas the other dependencies are caused indirectly by influencing intermediate variables. The hope that real world data are well explained by log-linear models can be backed up by the idea that each variable is influenced only by a few other variables directly. The graph of the variables with edges representing the direct influences is simple. This idea can be formulated using Markov networks, see Section 2.3.

### 2.2 Model Selection and AIC

A simple way of choosing the degree of the interaction terms in a log-linear model goes as follows. We begin with the simplest model class and decide on the basis of some significance test whether the data suggest to reject the model or not. If so then we include two-variable interactions. The same significance test is applied to the new model. The procedure ends if frequencies observed in the data and the probabilities given by the model do not differ significantly according to the significance test. Let us describe two popular significance tests (Christensen 1997, Goodman 1978).

Let $`j\{1,\mathrm{},k\}`$ be the probability space consisting of $`k`$ possible events. Let $`p_1,\mathrm{},p_k`$ be the probabilities of the model and $`m_1,\mathrm{},m_k`$ the observed frequencies. The sum $`l:=_jm_j`$ is the sample size. Then we calculate

1. Pearson chi-square Test:

$$X^2:=\underset{j}{}\frac{(m_jlp_j)^2}{lp_j}$$
2. Deviance Test:

$$G^2:=2\underset{j}{}m_j\mathrm{ln}\frac{m_j}{lp_j}.$$

Remark: $`G^2/(2l)`$ is the Kullback-Leibler distance (Cover and Thomas 1991) between the probability distribution $`(p_j)_{jk}`$ and the relative frequencies $`(m_j/l)_{jk}`$. Furthermore the minimizer of $`G^2`$ and $`R_{emp}`$ are the same, i.e., $`G^2/(2l)R_{emp}`$ does not depend on $`(p_j)_{jk}`$.

Under the assumption that $`(p_j)_{jk}`$ is the true probability distribution the values of the random variables $`X^2`$ and $`G^2`$ are in the large sample limit $`\chi ^2`$-distributed (Christensen 1997). Using this argument the model is rejected if the $`X^2`$ or $`G^2`$ values are too large. The range of acceptable values for $`X^2`$ and $`G^2`$ is based on the following rule:

First we have to calculate the number of degrees of freedom ($`df`$) of the model. ($`df`$) is given by the number of possible events minus $`1`$ minus the number of free parameters in the model class. Given $`df`$ one can look up in a table with $`\chi ^2`$ distribution whether the probability that $`X^2`$ or $`G^2`$, respectively, is outside a certain confidence interval.

However, the idea of rejecting a simple model only if the data contradicts the externally given significance level leads to conservative models. Suppose the outcome of the chi-squared test of a model yields values of $`\chi ^2`$ or $`G^2`$ within a certain confidence interval. Nevertheless there might be a model with only one interaction term more such that the $`\chi ^2`$ or $`G^2`$ value is considerably decreased. In this case we will certainly prefer the model that is a little bit more complex. This example illustrates that the externally given significance level implicitly determines the model selection. Therefore it is crucial to find a more systematic way for model selection and prevent in this way over- and underfitting.

Let $`(P_\theta )_{\theta ^k}`$ be a set of probability distributions with $`k`$ parameters and assume that the true probability distribution $`P`$ lies in $`(P_\theta )_{\theta ^k}`$. The true probability distribution $`P`$ is the minimizer of $`R`$ and for a sample size which is large compared to the size of $`(P_\theta )_{\theta ^k}`$, the empirical risk functional $`R_{emp}`$ is a good approximation for $`R`$. In this case we can search for the minimizer of $`R_{emp}`$ instead. This leads to the Deviance test above. However we are interested in the case where we cannot assure that the data set is sufficiently large. Then the minimizer of $`R_{emp}`$ is in general not a good estimation for the minimizer of $`R`$.

Akaike (1973) suggested to modify the empirical risk functional according to the number of free parameters of the model under consideration. He showed under some regularity conditions on the mapping $`\theta P_\theta `$, that

$$R_{emp}+df/l$$

is a consistent estimator of $`R`$, i.e., it has the same expectation value as $`R`$. However minimizing $`R_{emp}+df/l`$ for a finite data set give us no information about the true risk $`R`$. In Section 3.1 we obtain some bounds on the true risk and suggest to minimize this guaranteed risk rather than the expected risk, see Section 3.1.

### 2.3 Why Log-linear models

In Section 2 we have introduced log-linear model classes without giving any justification for it. In this section we explain why one should expect that probability distributions generating real world data can be described by $`k`$-factor models with small $`k`$. The key idea is that the distribution of a random variable $`X_j`$ is only influenced by the values a few others variables. This kind of statistical relations can be encoded by a Markov network.

The presentation of Markov networks in this work follows Pearl (1988). For each random variable $`X_i`$ we define the Markov boundary $`_i\{X_1,\mathrm{},X_{i1},X_{i+1},\mathrm{},X_n\}`$ as the smallest set fulfilling

$$P(x_i|x_1,\mathrm{},x_{i1},x_{i+1},\mathrm{},x_n)=P(x_i|_i).$$
(3)
Observe that this definition does not take any (temporal or causal) order of the random variables into account. It will therefore lead to an undirected graph. This is in contrast to the directed graphs in the context of Bayesian networks (Pearl 1985) which are for instance useful to formalize causal structure (Pearl 2000).

In case of strict positive probability distributions the Markov boundary is unique. Actually, strict positivity of the probability distributions is also a crucial condition for our approach to learn probabilities as we will see in Section 3.1. We now define the Markov network as the graph $`𝒢`$ with $`n`$ nodes also denoted by $`X_1,\mathrm{},X_n`$ and undirected edges between $`X_i`$ and $`X_j`$ if and only if $`X_jB_i`$ (or equivalently $`X_iB_j`$). In other words, we can generate the Markov network for a distribution $`P`$ if we connect each of the $`n`$ nodes with the element of its Markov boundary. As a consequence we obtain that two random variables $`X_i`$ and $`X_j`$, $`ij`$ are conditionally independent with respect to a set $`𝐙𝐗\{X_i,X_j\}`$, i.e.

$$P(X_i|X_j,Z)=P(X_i|Z),$$

if and only if every path from $`X_i`$ to $`X_j`$ has one node in $`𝐙`$. In other words, if $`𝐙`$ blocks every path between $`X_i`$ and $`X_j`$ then $`X_i`$ and $`X_j`$ are conditionally independent (with respect to $`𝐙`$). A clique in the graph $`𝒢`$ is a subgraph in which all nodes are connected to each other. We have an (partial) order on the cliques given by the inclusion as sets.

Let $`𝒢`$ be a Markov graph for a strictly positive probability distribution $`P`$ and let $`C`$ be the set of maximal cliques in $`𝒢`$. Then one can show that $`P`$ factorizes as follows

$$P(𝐱)=\frac{1}{N}\underset{cC}{}\varphi _c(𝐱_c),$$
(4)
where $`𝐗_c𝐗`$ is the set of random variables in the cliques $`c`$. $`N`$ is the normalization constant and $`\varphi _c`$ are the compatibility functions w.r.t. the clique $`c`$, also called the potential functions.

Suppose two distributions have the same Markov graph, then they describe the same (in)dependency in $`𝐗`$. In Section 3.1 we will see that the maximal number in a clique $`\mathrm{deg}(𝒢)=\mathrm{max}_{cC}\mathrm{\#}c`$ are a appropriate parameters to subdivide the set of all distributions $`𝒫(𝐗)`$ into  small model classes.

Consider the case that every clique has less than $`k+1`$ elements. Then the distribution can be described by a $`k`$-factor model due to the factorization in eq. (4). The Markov boundary can be understood in a rather literary sense if the variables represent some quantities that are measured at different positions in the real space. Then the belief to obtain probability distributions that correspond to simple Markov networks stems from the locality principle: distant variables influence each other only indirectly via other variables.

In statistical physics log-linear models can be justified even more directly. Assume the variable $`X_j`$ describes the physical state of particle $`j`$ and that the total energy of the system in the state $`𝐱=(x_1,\mathrm{},x_n)`$ is a sum of functions depending on at most $`k`$ particles, i.e.,

$$E(𝐱)=\underset{𝐣}{}f_𝐣(x_{j_1},\mathrm{},x_{j_k}),$$

where $`𝐣`$ runs over all $`k`$-subsets in $`\{1,\mathrm{},n\}`$. Then the thermodynamic equilibrium state, the so-called Gibbs distribution, is up to a normalization factor directly given by the exponential of the negative energy function. Hence the equilibrium distribution is exactly described by a $`k`$-factor model. In statistical physics one usually considers two particle interactions only, therefore $`2`$-factor models describe the equilibrium state.

Now we have justified why the class structure of log-linear models is natural. However, we have to refine our hierarchy in some respect. So far, we did not formulate any restrictions on the range of the probability distribution of the models. Intuitively it is clear that we can assign a small probability to an event $`𝐱\mathrm{\Omega }`$ only if we have a very large data set. This means that reasonable model selection should allow small probabilities only for large data sets. However, the empirical risk functional $`R_{emp}`$ penalizes for small probabilities only if they are contained in the data set. Therefore the value of $`R_{emp}`$ depends strongly on the specific data set if the probability distribution is not bounded from below by a strict positive constant. Actually, for deriving a bound for the risk functional $`R`$ in Section 3.1 we precisely need to bound the probability in this way.

###### Definition 1

Let $`𝒫_\lambda `$ be the set of all strict positive functions greater than $`\lambda >0`$. Let $`𝒫^k`$ be the set of probability distributions $`P`$ corresponding to a $`k`$-factor model, i.e.,

$$P(x_1,x_2,\mathrm{},x_n)=\underset{𝐣}{}q_𝐣(x_{j_1},x_{j_2},\mathrm{},x_{j_s}),$$

where each $`q_𝐣`$ is a positive function and $`𝐣:=(j_1,j_2,\mathrm{},j_s)`$ runs over all possible subsets of $`\{1,2,\mathrm{},n\}`$ with $`s`$ elements. Then we define the model class of degree $`k`$ by

$$_k^\lambda =𝒫_\lambda (𝐗)𝒫^k.$$
(5)
Remark: The values $`q_𝐣(x_{j_1},x_{j_2},\mathrm{},x_{j_s})`$ cannot be interpreted as probabilities since they can be greater than $`1`$.

the

## 3 LEARNING LOG-LINEAR <br>MODELS

It is a well-known problem in any mathematical theory of learning that a model can explain the training data well but does not fit the data observed in the future if the model class is too large. If the model class is small enough, one knows with high confidence that the model which minimizes the error on the training data will also explain the future data almost optimal within this model set.

### 3.1 Risk estimation

For a set of real-valued functions on an arbitrary set $`\mathrm{\Omega }`$ a relevant measure for the size of this set is the so-called VC dimension (Vapnik-Chervonenkis). First, let us define the VC dimension (Vapnik 1998) of a set of indicator functions $`(f_\alpha )`$, i.e. $`f_\alpha :\mathrm{\Omega }\{0,1\}`$.

###### Definition 2

Let $`\mathrm{\Lambda }`$ be an index set of arbitrary cardinality. Let $`(f_\alpha )_{\alpha \mathrm{\Lambda }}`$ be a set of indicator functions on $`\mathrm{\Omega }`$. Then the $`VC`$ dimension of $`(f_\alpha )_{\alpha \mathrm{\Lambda }}`$ is the largest number $`h`$ such that there exist $`h`$ points $`x_1,x_2,\mathrm{},x_h\mathrm{\Omega }`$ with the property that for every function
$`\chi :\{x_1,\mathrm{},x_h\}\{0,1\}`$ there exists a function $`f_\alpha `$ such that its restriction to $`\{x_1,x_2,\mathrm{},x_h\}`$ coincides with $`\chi `$.

The VC dimension for a set of real-valued functions is defined by the VC dimension for the set of corresponding indicator functions.

###### Definition 3

Let $`(f_\alpha )_{\alpha \mathrm{\Lambda }}`$ be a family of real-valued functions on a set $`\mathrm{\Omega }`$. Then the VC dimension of $`(f_\alpha )_{\alpha \mathrm{\Lambda }}`$ is the VC dimension of the family of the indicator functions $`(\chi _\mu f_\alpha )_{\mu ,\alpha \mathrm{\Lambda }}`$, where $`\chi _\mu `$ is $`\chi _\mu (x)=0`$ for $`x<\mu `$ and $`\chi _\mu (x)=1`$ for $`x\mu `$.

The notion of VC dimension plays a crucial role in statistical learning theory. It is a measure which indicates whether a learning machines overfits or not. Roughly speaking, if a family of functions which can be implemented by a learning machine has small VC dimension then every function which fits well the training data will fit with high probability the test data as well. The following theorem is an Corollary of Theorem 5.1 pp.192 in (Vapnik 1998).

###### Theorem 4

Let $`(f_\alpha )_{\alpha \mathrm{\Lambda }}`$ be a measurable set of bounded real-valued functions on $`\mathrm{\Omega }`$, $`Af_\alpha (z)B`$, and let the set of indicator functions have finite VC dimension $`h`$. Let $`\mu `$ be a probability distribution on $`\mathrm{\Omega }\times `$ and let the data $`𝐱_1,𝐱_2,\mathrm{},𝐱_l`$ be drawn according to $`\mu `$ (independent and identically distributed). Then we have with probability at least $`1\eta `$

$$\stackrel{~}{R}(f_\alpha )\stackrel{~}{R}_{emp}(f_\alpha )+(BA)\sqrt{\frac{h\mathrm{ln}h+\mathrm{ln}16l\mathrm{ln}\eta }{l}},$$

where $`\stackrel{~}{R}_{emp}(f_\alpha ):=\frac{1}{l}_{jl}f_\alpha (𝐱_𝐣)`$ is the empirical risk and $`\stackrel{~}{R}(f_\alpha ):=f_\alpha (𝐱)𝑑\mu (𝐱)`$ the true risk.

Remark: We have $`R(P_\alpha )=\stackrel{~}{R}(\mathrm{ln}P_\alpha )`$ and $`R_{emp}(P_\alpha )=\stackrel{~}{R}_{emp}(\mathrm{ln}P_\alpha )`$ when $`(P_\alpha )`$ is a set of probability distributions. Therefore we can use this theorem to calculate the generalization error of equation (1) for the family of probability distributions defined in equation (5).

First, let us interpret the generalization error in this context. Suppose we have found a probability distribution $`P_t`$ in some model class that minimizes $`R`$ then we know that in this space $`P_t`$ is the closest distribution to the true distribution in the Kullback-Leibler distance. Furthermore, we can conclude with high confidence that future data will fit to this distribution as well.

Next we estimate the VC dimension of product distributions since this will be the leading intuition for estimating the VC dimension of the model class $`_k^\lambda `$.

###### Lemma 1

Let $`X_1,X_2,\mathrm{},X_n`$ be discrete random variables where $`X_i`$ takes $`m_i`$ different values in $`\mathrm{\Omega }_i`$ for $`i=1,\mathrm{},n`$. Then the VC dimension of the family of all product distributions $`𝒫^1=\{P(x_1,\mathrm{},x_n)=P(x_1)\mathrm{}P(x_n):P\text{ is a probability distribution}\}`$ is given by $`N=1+_{i=1}^nm_in`$.

Remark: For the families $`(\mathrm{ln}P)_{P𝒫^1}`$ and $`(P)_{P𝒫^1}`$ the set of indicator functions coincides. Hence $`(\mathrm{ln}P)_{P𝒫^1}`$ and $`(P)_{P𝒫^1}`$ have the same VC dimension.

Proof: Let $`x_{j;0}`$ be an arbitrary element of the set $`\mathrm{\Omega }_j`$. Due to $`P(𝐱)=P(x_1)\mathrm{}P(x_n)`$ we can write the logarithm of the joint probability as

$$\mathrm{ln}P(𝐱)=\underset{j}{}\mathrm{ln}(P(x_j)/P(x_{j;0}))+\underset{j}{}\mathrm{ln}P(x_{j;0})$$

We can characterize an $`n`$-tuple $`𝐱=(x_1,x_2,\mathrm{},x_n)`$ uniquely by an $`N=_j(m_j1)`$ dimensional vector with entries $`0`$ and $`1`$ as follows: The vector consists of $`n`$ blocks of dimension $`(m_j1)`$ characterizing $`x_j`$. Let $`x_{j;0},x_{j:1},\mathrm{},x_{j;m_j1}`$ be the elements of $`\mathrm{\Omega }_j`$ in an arbitrary ordering. For $`x_j=x_{j;i}`$ with $`i0`$ let the $`i`$-th entry of the $`j`$-th block be $`1`$ and the other entries of the $`j`$-th block be zero. For $`x_j=x_{j;0}`$ let all entries of the $`j`$-th block be zero.

For each $`n`$-tuple the logarithm of its probability is given as follows. Let $`c^𝐱^N`$ be the vector corresponding to $`𝐱:=(x_1,x_2,\mathrm{},x_n)`$. Then we have

$$\mathrm{ln}P(𝐱)=\underset{l}{}c_l^𝐱f_l+d=c^𝐱,f+d$$

where the coordinates $`f_l`$ of $`f`$ are given by the values $`\mathrm{ln}(P(x_{j;i})/P(x_{j;0}))`$ with $`jn`$ and $`1im_j1`$. The constant term $`d`$ is defined as $`d:=_j\mathrm{ln}P(x_{j;0})`$. This shows that $`\mathrm{ln}P(𝐱)`$ can be written as an affine functional in $`^N`$. Hence the VC dimension of $`𝒫`$ is at most $`N+1`$ (see (Vapnik 1998), Example in Chapter 5.2.3). To show that it is not smaller than $`N+1`$ note that there is no restriction at all to the set of all possible vectors $`f`$ if $`d`$ is chosen appropriately. This can be seen as follows. Let $`g_1,g_2,\mathrm{},g_{m_j1}`$ be the entries of $`f`$ in the $`j`$-th block. Then set

$$P(x_{j;i}):=\frac{\mathrm{exp}(g_i)}{_{1im_j1}\mathrm{exp}(g_i)+1}$$

and

$$P(x_{j;0}):=\frac{1}{_{1im1}\mathrm{exp}(g_i)+1}.$$

Now we have to show that one can find $`N+1`$ vectors that can be classified in all $`2^{N+1}`$ possible ways such that the vectors correspond to $`n`$-tuples. Consider the vectors $`e_j`$ having $`1`$ at the position $`j`$ and $`0`$ elsewhere. Obviously, each vector $`e_j`$ corresponds to a possible $`n`$-tuple. As $`N+1`$-th vector we choose the origin $`(0,0,\mathrm{},0)`$. It corresponds to the $`n`$-tuple $`(x_{1;0},x_{2;0},\mathrm{},x_{n;0})`$. These vectors can be classified in all $`2^{N+1}`$ possibilities as follows. Choose an arbitrary indicator function $`\chi `$ on these $`N+1`$ points and set $`f:=(\pm 1,\pm 1,\mathrm{},\pm 1)`$ with positive sign at position $`j`$ if and only if $`\chi (e_j)=1`$. Define $`d`$ as above in order to ensure that

$$𝐱c^𝐱,f+d$$

defines the logarithm of a probability function. Set $`a:=d\pm 1/2`$ and choose the sign positive if and only if $`\chi ((0,0,\mathrm{},0))=1`$. Then the indicator function $`\chi _aw`$ with $`w(𝐱):=c^𝐱,f+d`$ coincides with the desired classifier $`\chi `$. Hence the VC dimension of product distributions is $`N+1`$. $`\mathrm{}`$

Now we give a bound for the VC dimension of the model classes introduced in Section 2.3.

###### Lemma 2

The VC dimension of $`𝒫^k`$ is at most

$$h_k:=\underset{𝐣}{}m_{j_1}m_{j_2}\mathrm{}m_{j_k},$$

where $`m_i`$ is the size of the set $`\mathrm{\Omega }_i`$.

Remark: $`h_s`$ is at most $`O(n^k)`$ if $`m_i`$ is uniformly bounded.

Proof: Let $`P𝒫^k`$. Since $`P`$ factorizes we can characterize the expression $`\mathrm{ln}P`$ uniquely by a vector $`f^h`$ as follows: the vector $`f`$ consists of $`\left(\genfrac{}{}{0pt}{}{n}{k}\right)`$ blocks where the blocks are indexed by $`𝐣`$ given in the Lemma. The block $`𝐣`$ has size $`_im_{j_i}`$ and the entries within the block have the values

$$\mathrm{ln}q_𝐣(x_{j_1},x_{j_2},\mathrm{},x_{j_k}),$$

where $`(x_{j_1},x_{j_2},\mathrm{},x_{j_k})`$ runs over all elements of $`\mathrm{\Omega }_{j_1}\times \mathrm{\Omega }_{j_2}\times \mathrm{}\times \mathrm{\Omega }_{j_k}`$. Each $`n`$-tuple $`𝐱=(x_1,x_2,\mathrm{},x_n)`$ can be characterized by a vector $`c^𝐱^h`$ as follows: In each block $`𝐣`$ there is exact one coordinate that corresponds to the $`s`$-tuple $`(x_{j_1},x_{j_2},\mathrm{},x_{j_k})`$. Set this entry to $`1`$ and the other entries of the block to $`0`$. Then $`\mathrm{ln}P(𝐱)`$ is given by the inner product of the vectors $`f`$ and $`c^𝐱`$. We obtain the lemma since the VC dimension of the set of linear functionals on $`^h`$ is $`h`$, see (Vapnik 1998). $`\mathrm{}`$

Obviously, the model class $`_k^\lambda `$ is a subset of $`𝒫^k`$ and so we can give an upper bound the VC dimension of $`_k^\lambda `$ with $`h_k`$ which is independent of $`\lambda `$. Therefore Lemma 2 and Theorem 4 yields to the following theorem.

###### Theorem 5

Let $`_k^\lambda `$ be the model class of degree $`k`$ and $`h_k`$ be the upper bound on the VC dimension of $`_k^\lambda `$ given in Lemma 2. Then for a training data $`𝐱_1,𝐱_2,\mathrm{},𝐱_l`$ we have for $`P_\alpha _k^\lambda `$ with probability at least $`1\eta `$

$$R(P_\alpha )R_{emp}(P_\alpha )+\varphi (k,\lambda ,\eta ),$$
(6)
with $`\varphi (k,\lambda ,\eta ):=\mathrm{ln}\lambda \sqrt{\frac{h_k\mathrm{ln}h_k+\mathrm{ln}16+\mathrm{ln}l\mathrm{ln}\eta }{l}}`$.

The task to find the minimizer of $`R_{emp}`$ in $`_k^\lambda `$ even for fixed $`\lambda `$ and $`k`$ is computationally expensive. We will not treat this question here. Theorem 5 tells us how to estimate the risk $`R(P_\alpha )`$ for a a priori chosen model class $`_k^\lambda `$. However we do not know which class to select. This question is addressed in the next section.

### 3.2 Estimating Log-linear Models from empirical data

We do not want to exclude $`k`$-factor models with large $`k`$ a priori if the data strongly indicates that such a complex model is appropriate. Similarly, we do not want to exclude distributions with small $`\lambda `$. We rather want to avoid such models if the data does not give us strong enough reason for choosing them. To solve this problem we use a generalized version of the structural risk minimization principle (Vapnik 1998, Vapnik 1995). In this way the result of Theorem 5 can be extended such that $`k`$ and $`\lambda `$ need not to be chosen in advanced. For doing so we have to define a prior probability measure $`\nu `$ on the set of possible pairs $`(k,\lambda )`$ in advance in order to get a bound on the true risk.

We obtain from Theorem 5 with a standard union bound argument, see e.g. Lemma 4.1 in (Herbrich 2002), following corollary.

###### Corollary 1

Let $`(\lambda _n)_n`$ be a strict monotone decreasing sequence converging to zero. Let $`\nu `$ be a strict positive probability measure on $`^2`$ and $`A_{kn}=\{1,\mathrm{}k\}\times \{1,\mathrm{}n\}`$. Then for training data $`𝐱_1,𝐱_2,\mathrm{},𝐱_l`$ we have with probability at least $`1\eta `$

$$R(P)R_{emp}(P)+\varphi (k,\lambda _n,\eta \nu (A_{kn})),$$

for $`P_k^{\lambda _n}`$ and $`\varphi `$ defined above.

The probability measure $`\nu `$ in Corollary 1 has to be chosen before seeing the training data. Usually one chooses $`\nu `$ to be inversely proportional to $`\varphi `$. However, it should be emphasized that the statement in Corollary 1 does not assume that the measure $`\nu `$ is chosen according to any prior probabilities for the possible probability measures as used in Bayesian learning.

Now we can formulate structural risk minimization as follows:

$$\underset{(k,n)^2}{\text{Minimize}}t+\varphi (k,\lambda ,\eta \nu (A_{kn}))$$

$$\text{subject to }t=\underset{P_k^{\lambda _n}}{\mathrm{min}}R_{emp}(P)$$

### 3.3 Model Computation by Convex Optimization

In this section we show that the minimization of the empirical risk $`R_{emp}`$ is a convex optimization problem in the class of all $`k`$-factor models for a fixed $`k`$ and $`\lambda `$. We formulate the optimization problem for $`\mathrm{ln}P`$ rather than for $`P`$. Let $`𝐱^1,\mathrm{},𝐱^l`$ be the training data. For each set $`𝐣=\{j_1,\mathrm{},j_k\}\{1,\mathrm{},n\}`$ let $`V_𝐣`$ be the vector space of real-valued functions on the set $`\mathrm{\Omega }_𝐣:=\mathrm{\Omega }_{j_1}\times \mathrm{}\times \mathrm{\Omega }_{j_k}.`$ Define

$$V^k:=_𝐣V_𝐣$$

where $`𝐣`$ runs over all subsets of $`\{1,\mathrm{},n\}`$ of size $`k`$. Now fix $`k`$ and $`𝐣`$ of size $`k`$. The set $`\mathrm{\Omega }_𝐣`$ is discrete and so the functions

$$e_{w_1,w_2,\mathrm{},w_k}(v_1,v_2,\mathrm{},v_k)=\{\begin{array}{cc}1\hfill & \text{if }w_1=v_1,\mathrm{},w_k=v_k\hfill \\ 0\hfill & \text{else}\hfill \end{array}$$

form a canonical basis of $`V_𝐣`$. Further let $`𝐱_𝐣`$ be the $`k`$-tuple given by the restriction of $`𝐱\mathrm{\Omega }`$ to the variables $`x_{j_1},\mathrm{},x_{j_k}`$ and $`𝐱_𝐣^i`$ the restriction of the training data $`𝐱^i`$ to the variables of $`𝐱_𝐣`$.

Now, let us formulate the optimization problem of the empirical risk. Find a vector $`f=_𝐣f_𝐣V^k`$ with $`f_𝐣V_𝐣`$ such that

$$\stackrel{~}{R}_{emp}(f):=\frac{1}{l}\underset{i=1}{\overset{l}{}}\underset{𝐣}{}f_𝐣(𝐱_𝐣^i)$$
(7)
is minimal subject to the constraints

$$Z(f):=\underset{𝐱𝛀}{}\mathrm{exp}(\underset{𝐣}{}f_𝐣(𝐱_𝐣))=1$$

and

$$G_𝐱(f):=\underset{𝐣}{}f_l(𝐱_𝐣)\mathrm{log}\lambda \text{for all }𝐱\mathrm{\Omega }\text{.}$$

Then the probability measure $`P_k^\lambda `$ is given by $`\mathrm{ln}f`$. The functions $`\stackrel{~}{R}_{emp}`$ and $`G_𝐱`$ are obviously linear in $`f`$. The function $`Z`$ is convex in $`f`$ since the exponent is linear in $`f`$ and the exponential function is convex. Clearly, the sum of the convex functions

$$f\mathrm{exp}(\underset{𝐣}{}f_𝐣(𝐱_𝐣))$$

over all $`𝐱`$ is convex. Therefore the problem can be solved with standard convex optimization methods (Pallaschke and Rolewicz 1997). However, for a large number of variables the solution becomes computationally expensive. Notice that the number of constraints in (7) increases exponentially with $`n`$. Furthermore the calculation of $`\stackrel{~}{R}_{emp}(f)`$ involves the summation over exponentially many possible $`n`$-tuples $`𝐱`$. This problem is similar to the calculation of the partition function in statistical mechanics which requires a summation over exponentially many states of a many-particle system. Those computational problems are not different from those appearing in conventional approaches to log-linear model selection.

## 4 CONCLUSION

We used structural risk minimization of statistical learning theory to obtain a new selection criterion for log-linear models. In this we way have not only a criterion to choose the suitable model complexity, but we found also a bound on the actual risk, i.e. the key feature of our approach is that it provides statements of the form “if a log-linear model of certain simplicity fits well to the observed data then we are guaranteed with high probability that it will fit future observations as well”. This kind of nonasymptotic statistics become more and more important for in research areas like cognitive and social science, where the number of possibly relevant features for the system under investigation is large. It offers a new way of developing models, namely, infering statistical dependencies amoung a large number of features from data.

We structured the class of log-linear models according to the degree of the interaction terms. Using other measures of complexity than VC dimension like Gaussian complexity (Bartlett and Mendelson 2001) we believe that the obtained bounds can be improved. The structure of the model classes can be refined taking for example prior knowledge into account. In future this model selection criterion need to be tested on real-world data and the results compared to existing ones.


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# The 𝜂_𝑐(3654) and hyperfine splitting in charmonium

## I Introduction

Recently the Belle Collaboration has observed a new charmonium state, the $`\eta _c(2S)`$, in exclusive $`BKK_SK^{}\pi ^+`$ decays ref.01 . The measured mass of the $`\eta _c(2S)`$, $`M(2^1S_0)=3654\pm 14`$ MeV, has appeared to be rather close to the $`\psi (2S)`$ mass, $`M_2=3686`$ MeV, giving rise to a small hyperfine splitting for the $`2S`$ state,

$$\mathrm{\Delta }_{\mathrm{HF}}(2S,\mathrm{exp})=M_2M(2^1S_0)=32\pm 14\mathrm{MeV}.$$
(1)
Compared to the hyperfine splitting for the $`1S`$ state, $`\mathrm{\Delta }_{\mathrm{HF}}(1S,\mathrm{exp})=117.2\pm 1.5`$ MeV ref.02 , the number (1) is 3.5 times smaller,

$$\frac{\mathrm{\Delta }_{\mathrm{HF}}(2S,\mathrm{exp})}{\mathrm{\Delta }_{\mathrm{HF}}(1S,\mathrm{exp})}=0.273\pm 0.123$$
(2)
and has a large experimental error coming from the error in the $`\eta _c(2S)`$ mass. However, in the ratio (2) even the upper limit is rather small and it is of interest to compare this small ratio with the predictions in the conventional theory of the spin-spin interaction in QCD.

In this rapid communication we shall discuss the problem posed in a model-independent way and show that the experimental value (2) puts a strong restriction on the coupling $`\alpha _{\mathrm{HF}}(2S)`$, determining the spin-spin interaction. In particular, the central value of $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ corresponds to $`\alpha _{\mathrm{HF}}(2S)0.18÷0.19`$ which implies the very large value for the renormalization scale $`\mu _2=2M_27.4`$ GeV. This scale appears to be drastically different from that for the $`1S`$ state $`(J/\psi )`$ where the scale in the strong coupling determining the HF interaction is $`\mu _1\frac{1}{2}M_1=1.55`$ GeV $`(M_1=M(J/\psi ))`$ and the coupling constant $`\alpha _{\mathrm{HF}}(\mu _1)0.335`$ is rather large. We shall show that if one uses for the $`2S`$ state the same prescription as for $`J/\psi `$, then our theoretical prediction is $`\mathrm{\Delta }_{\mathrm{HF}}(2S,\mathrm{theory})=57\pm 8`$ MeV, which by derivation is the lower bound of $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ and gives rise to the mass value $`M(\eta _c(2S))=3630\pm 8`$ MeV, being two standard deviations off the value obtained from the Belle experiment ref.01 .

## II Model-independent calculation of $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$

In one-loop approximation in the $`\overline{MS}`$ renormalization scheme, the hyperfine splitting, $`\mathrm{\Delta }_{\mathrm{HF}}(nS)`$, is given by the well-known expression:

$$\mathrm{\Delta }_{\mathrm{HF}}(nS)=\frac{8}{9}\frac{\alpha _{\mathrm{HF}}(\mu _n)}{m_c^2}|R_n(0)|^2\left(1+\frac{\alpha _{\mathrm{HF}}(\mu _n)}{\pi }\xi _{\mathrm{HF}}\right),$$
(3)
where the factor $`\xi _{\mathrm{HF}}`$ comes from the one-loop corrections, $`\xi _{\mathrm{HF}}=\frac{5}{12}\beta _0\frac{8}{3}\frac{3}{4}\mathrm{ln}2`$ ref.03 , and for $`n_f=4`$, $`\xi _{\mathrm{HF}}(n_f=4)=0.2857`$ is small, so that the $`\alpha _{\mathrm{HF}}`$ correction in the brackets turns out to be less than 3% $`(\alpha _{\mathrm{HF}}(\mu _n)\begin{array}{c}<\hfill \\ \hfill \end{array}0.35)`$ and can be neglected in the ratio (2).

For light mesons the relativistic version of the expression (3) also exists, it was derived in Refs. ref.04  and can be useful for charmonium, since the current $`c`$-quark mass is not large, $`m_c=1.3.\pm 0.2`$ GeV ref.02 ,

$$\mathrm{\Delta }_{\mathrm{HF}}^R(nS)=\frac{8}{9}\frac{\alpha _{\mathrm{HF}}(\mu _n)}{\omega _n^2}|R_n(0)|^2.$$
(4)
Here $`\omega _n`$ is the kinetic energy matrix element:

$$\omega _n=\sqrt{p^2+m_c^2}_{nS},$$
(5)
which plays the role of the constituent quark mass. The HF splitting (3) or (4) strongly depends on the $`c`$-quark mass chosen and $`\omega _n^2`$ can change by a factor of 2 for different choices of $`m_c`$. To escape the problem of a correct choice of $`m_c`$ it is convenient to consider the ratio of hyperfine splittings with relativistic corrections:

$$\frac{\mathrm{\Delta }_{\mathrm{HF}}^R(2S)}{\mathrm{\Delta }_{\mathrm{HF}}^R(1S)}=\frac{\alpha _{\mathrm{HF}}(\mu _2)}{\alpha _{\mathrm{HF}}(\mu _1)}\left(\frac{\omega _1}{\omega _2}\right)^2\left|\frac{R_2(0)}{R_1(0)}\right|^2.$$
(6)
Note that the difference between $`\omega _1`$ and $`\omega _2`$ in charmonium which takes into account relativistic corrections, is not large, e.g., for the spinless Salpeter equation $`(\omega _1/\omega _2)^20.94`$, however, it will be useful to keep this factor in our analysis, since due to this factor the ratio (6) and therefore the hyperfine splitting is 6% smaller. However, the difference between $`\omega (nS)`$ and the current (pole) mass $`m_c`$ coming from relativistic corrections, is not small, being about 200 MeV and 250 MeV for the 1S and 2S states, respectively (see Section IV).

The ratio of the wave functions at the origin occurring in Eq. (6) can be extracted in a model-independent way from the leptonic widths:

$$\mathrm{\Gamma }_{e^+e^{}}(nS)=\frac{4\alpha ^2e_q^2}{M_n^2}|R_n(0)|^2\gamma _n,$$
(7)
where the factor $`\gamma _n`$ is

$$\gamma _n=1\frac{16}{3\pi }\alpha _s(M_n),$$
(8)
and in $`\gamma _n`$ the strong coupling $`\alpha _s`$ is taken at the scale $`\mu _n=M_n`$ ($`M_1=M(J/\psi )`$, $`M_2=M(\psi (2S))`$. Note that just this prescription for $`\alpha _s(\mu )`$ (in the QCD factor $`\gamma `$) provides the minimal value of the extracted wave function at the origin. For $`\mathrm{\Lambda }^{(4)}`$(3-loop) $`=280`$ MeV ($`n_f=4`$), which corresponds to $`\alpha _s(M_z)=0.117`$, the values of $`\alpha _s(M_n)`$ in 3-loop approximation are the following,

$`\alpha _s(M_1)`$ $`=`$ $`0.247,\alpha _s(M_2)=0.232,`$
$`\gamma _1`$ $`=`$ $`0.581,\gamma _2=0.606.`$ (9)
Note that for another choice: $`\gamma _1=\gamma _2`$ which is often used, the value of $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ would be 4% larger than in our consideration.

Then from the ratio of leptonic widths (7) it follows that

$$\left|\frac{R_2(0)}{R_1(0)}\right|^2=\left(\frac{M_2}{M_1}\right)^2\frac{\mathrm{\Gamma }_{e^+e^{}}(2S)}{\mathrm{\Gamma }_{e^+e^{}}(1S)}\frac{\gamma _1}{\gamma _2}.$$
(10)
Taking the experimental values: $`\mathrm{\Gamma }_{e^+e^{}}(1S)=5.26\pm 0.37`$ keV and $`\mathrm{\Gamma }_{e^+e^{}}(2S)=2.19\pm 0.15`$ keV ref.02 , this ratio with a good accuracy is

$$\left|\frac{R_2(0)}{R_1(0)}\right|^2=(0.59\pm 0.08)\frac{\gamma _1}{\gamma _2}.$$
(11)
The values of the wave functions at the origin extracted from the leptonic widths are the following: $`|R_1(0)|^2=0.917`$ GeV3 and $`|R_2(0)|^2=0.518`$ GeV3 and at this point it is important to stress that for the $`2S`$ state the extracted value of the wave function at the origin already takes into account the influence of the close lying $`D\overline{D}`$ channel on the wave function in an implicit way, and it is the same both for the hyperfine splitting and the leptonic width. Due to the nearby $`D\overline{D}`$ threshold the experimental value of $`|R(0)|^2`$ as well as the leptonic widths typically appear to be smaller by about 20% than in the theoretical calculations (one-channel approximation) and the influence of the $`D\overline{D}`$ channel on the $`\psi (2S)`$ mass was discussed in Ref. ref.05 .

Now combining the expressions (11) and (6) one obtains

$$\frac{\mathrm{\Delta }_{\mathrm{HF}}^R(2S)}{\mathrm{\Delta }_{\mathrm{HF}}^R(1S)}=(0.59\pm 0.08)\eta _{\mathrm{HF}}$$
(12)
with the factor

$$\eta _{\mathrm{HF}}=\left(\frac{\omega _1}{\omega _2}\right)^2\frac{\alpha _{\mathrm{HF}}(\mu _2)\gamma _1}{\alpha _{\mathrm{HF}}(\mu _1)\gamma _2}.$$
(13)
From the analysis of the leptonic width and hyperfine splitting for the $`J/\psi `$ it is known that the renormalization scale $`\mu _1`$ in $`\alpha _{\mathrm{HF}}(1S)`$ is different from $`M_1`$ and a good description of the $`\mathrm{\Delta }_{\mathrm{HF}}(1S)`$ can be reached if in Eq. (3) or (4) the scale is taken to be $`\mu _1\frac{1}{2}M_1=1.55`$ GeV, which with $`\mathrm{\Lambda }^{(4)}(3\mathrm{loop})=280`$ MeV gives the value $`\alpha _{\mathrm{HF}}(\frac{1}{2}M_1)=0.335`$. Now with the same prescription $`\mu _2=\frac{1}{2}M_2`$, the value $`\alpha _{\mathrm{HF}}(\frac{1}{2}M_2)=0.307`$ is obtained, and then from Eqs. (9) and (13) it follows that

$$\eta _{\mathrm{HF}}=0.826,\mathrm{for}\left(\frac{\omega _1}{\omega _2}\right)^2=0.94,\frac{\gamma _1}{\gamma _2}=0.96.$$
(14)
As a result from Eq. (12) the ”theoretical” ratio of the hyperfine splittings is

$$\frac{\mathrm{\Delta }_{\mathrm{HF}}^R(2S)}{\mathrm{\Delta }_{\mathrm{HF}}^R(1S)}=(0.48\pm 0.07)$$
(15)
and our prediction for the $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ is

$$\mathrm{\Delta }_{\mathrm{HF}}(2S,\mathrm{theory})=(57\pm 8)\mathrm{MeV}.$$
(16)
Owing to our derivation this number can be considered as the lower bound of $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ which was obtained in a model-independent way. The predicted mass of the $`\eta _c(2S)`$, $`M(\eta _c(2S))=3630\pm 8`$ MeV, differs by two standard deviations from the measured value in the Belle experiment ref.01 , but is rather close to the $`\eta _c(2S)`$ mass observed in another Belle experiment ref.06  $`e^+e^{}J/\psi \eta _c`$ where the measured $`\eta _c(2S)`$ mass, $`M(\eta _c(2S))=3622\pm 12`$ MeV is in good agreement with our prediction. Note that besides the dominant perturbative term in the HF splitting there exists also a nonperturbative (NP) contribution to $`\mathrm{\Delta }_{\mathrm{HF}}(nS)`$. However, the NP terms in charmonium can be calculated as in Ref. ref.08  and turn out to be small: we find $`\mathrm{\Delta }_{\mathrm{HF}}^{\mathrm{NP}}(1S)=35`$ MeV and $`\mathrm{\Delta }_{\mathrm{HF}}^{\mathrm{NP}}(2S)=12`$ MeV. It is of interest to note that the same model-independent estimate of the $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ only with $`\gamma _1=\gamma _2`$ and $`\alpha _{\mathrm{HF}}(\mu _1)=\alpha _{\mathrm{HF}}(\mu _2)`$ was suggested many years ago at the time when the mass of the $`\eta _c(1S)`$ was not correctly measured ref.07 .

## III The renormalization scale in $`\alpha _{\mathrm{HF}}(\mu )`$

To obtain the central value in Eq. (1) in the same procedure one needs to take the very small value $`\alpha _{\mathrm{HF}}(\mu _2)=0.187`$, which corresponds to the large renormalization scale: $`\mu _2=7.4`$ GeV $`\begin{array}{c}>\hfill \\ \hfill \end{array}2M_2`$. It is difficult to point out any physical explanation for such different scales of $`\mu _2=2M_2`$ for the 2S state and and $`\mu _1\frac{1}{2}M_1=1.55`$ GeV for the $`J/\psi `$.

Existing theoretical calculations of $`\mathrm{\Delta }_{\mathrm{HF}}`$ mostly provide the $`2S`$ splittings in the range 70-90 MeV ref.09 ; ref.10 , while only in Ref. ref.11 , where the modified (screened) color Coulomb interaction was used, the calculated HF splitting is smaller: $`\mathrm{\Delta }_{\mathrm{HF}}(2S)=38`$ MeV. However, in Ref. ref.11  with the same modified Coulomb interaction the fine structure splittings for the $`\chi _c(1P)`$ mesons turned out to be twice as small as the experimental data.

## IV Charmonium spectrum

Here we give as an illustration the spectrum and the constituent masses $`\omega _n`$ in charmonium which were calculated solving the spinless Salpeter equation with a linear plus Coulomb potential,

$$\left(2\sqrt{\widehat{p}^2+m_c^2}+V_0(r)\right)\psi _{nL}(r)=M_{nL}^{(0)}\psi _{nL}(r).$$
(17)
In the static potential

$$V_0(r)=\frac{4}{3}\frac{\alpha _{\mathrm{st}}}{r}+\sigma r$$
(18)
the parameters $`\alpha _{\mathrm{st}}=0.42`$ and $`\sigma =0.18`$ GeV2 were taken. For the one-loop pole mass we used $`m_c=1.42`$ GeV. The remarkable feature of this set of parameters is that the meson mass, $`M(nL)=M^{(0)}(nL)C_{SE}(nL)`$, contains a very small negative subtractive constant which can be strictly determined by the nonperturbative self-energy contribution: $`C_{SE}=4\sigma 0.24/\pi \omega _{nL}`$ ref.12  and its value $`C_{SE}=33\pm 3`$ MeV is small and approximately equal for all $`nS`$ states $`(n5)`$.

From Table 1 one can see that the energy $`\omega (nS)`$ of the $`c`$ quark, playing the role of the constituent mass of a given $`nS`$ state, appears to be around 1.65-1.75 GeV, i.e., essentially larger than the current mass, $`m_c1.4`$ GeV in the Salpeter equation and gives rise to a small suppression of the HF splittings. This fact can be considered as an explanation why in the nonrelativistic approach the constituent $`c`$-quark mass is usually taken to be larger, e.g., in Ref. ref.09  $`m_c=1.84`$ GeV. Note also that with $`\omega (1S)=1.60`$ GeV and neglecting the NP contribution of about 3-5 MeV, the value of $`\alpha _{\mathrm{HF}}(1S)=0.355`$ is needed in order to obtain the experimental number for $`\mathrm{\Delta }_{\mathrm{HF}}(1S)=117`$ MeV.

## V Conclusions

Thus we can conclude that the Belle experiment ref.01  on the $`\eta _c(2S)`$ mass if it is confirmed, creates a number of interesting theoretical problems, in particular, about the correct choice of the renormalization scale of $`\alpha _s(\mu )`$ for the excited states. In the conventional approach $`\alpha _{\mathrm{HF}}(\mu _2)0.30`$ and $`\mathrm{\Delta }_{\mathrm{HF}}(2S)(57\pm 8)`$ MeV and our prediction for the mass is $`M(\eta _c(2S))=3630\pm 88`$ MeV which appears to be in agreement with the result of the other Belle experiment ref.06 : $`e^+e^{}J/\psi ,\eta _c`$. To obtain decisive conclusions it would be important to measure the $`\eta _c(2S)`$ mass with better accuracy. The confirmation of the value of $`\mathrm{\Delta }_{\mathrm{HF}}(2S)`$ in the range close to 32 MeV would require a drastic reconsideration of our understanding of the scale in the HF interaction and possibly also of the QCD factor $`\gamma _n`$ present in the leptonic width, since they are interconnected.

###### Acknowledgements.

A.M.B. is grateful to L.B. Okun for his kind interest in this work and wants to thank both him and A. Drutskoy for useful discussions. This work was partly supported by the INTAS grant 00-00110 and the RFFI grant 00-02-17836.


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# Slopes of 2-adic overconvergent modular forms with small level

## 1 Abstract

Let $`\tau `$ be the primitive Dirichlet character of conductor $`4`$, let $`\chi `$ be the primitive even Dirichlet character of conductor $`8`$ and let $`k`$ be an integer. Then the $`U_2`$ operator acting on cuspidal overconvergent modular forms of weight $`2k+1`$ and character $`\tau `$ has slopes in the arithmetic progression $`\{2,4,\mathrm{},2n,\mathrm{}\}`$, and the $`U_2`$ operator acting on cuspidal overconvergent modular forms of weight $`k`$ and character $`\chi \tau ^k`$ has slopes in the arithmetic progression $`\{1,2,\mathrm{},n,\mathrm{}\}`$.

We then show that the characteristic polynomials of the Hecke operators $`U_2`$ and $`T_p`$ acting on the space of classical cusp forms of weight $`k`$ and character either $`\tau `$ or $`\chi \tau ^k`$ split completely over $`𝐐_2`$.

## 2 Introduction

###### Definition 1

Let $`f`$ be a cuspidal modular eigenform with $`q`$-expansion at $`\mathrm{}`$ given by $`_{n=1}^{\mathrm{}}a_nq^n`$. Let $`f`$ be *normalised*; that is, $`a_1=1`$. The *$`(p)`$-slope* of $`f`$ is defined to be the $`p`$-valuation of $`a_p`$; we normalise the $`p`$-valuation of $`p`$ to be $`1`$. If we do not specify $`p`$, then we mean the $`2`$-slope.

In this paper, we prove the following theorem on the slopes of classical modular cusp forms :

###### Theorem 2

Let $`\tau `$ be the nontrivial character of conductor $`4`$, and let $`k`$ be an integer greater than $`2`$. The slopes of the $`U_2`$ operator acting on $`S_{2k1}(\mathrm{\Gamma }_0(4),\tau )`$ are

$$2,4,6,\mathrm{},2k4.$$

Let $`\chi `$ be the even primitive Dirichlet character of conductor $`8`$. The slopes of the $`U_2`$ operator acting on $`S_k(\mathrm{\Gamma }_0(8),\chi \tau ^k)`$ are

$$1,2,3,\mathrm{},k2.$$

As a corollary of this theorem, we also prove the following result about the field over which cusp forms of weight $`k`$ and character $`\chi \tau ^k`$ or $`\tau `$ are defined:

###### Corollary 3

Let $`k`$ be an integer greater than $`2`$, and let $`S`$ be either $`S_{2k1}(\mathrm{\Gamma }_0(4),\tau )`$ or $`S_k(\mathrm{\Gamma }_0(8),\chi \tau ^k)`$

Let $`fS`$ be a normalised eigenform. Then the coefficients of the Fourier expansion of $`f`$ are elements of $`𝐐_2`$.

This corollary gives a partial answer to an extension of Questions 4.3 and 4.4 of Buzzard , which give a conjectural bound on the degree of the field of definition of certain spaces of modular forms over $`𝐐_2`$.

## 3 Previous work

Matthew Emerton determines in his thesis  the smallest slope for the spaces of modular cuspforms $`S_k(\mathrm{\Gamma }_0(2^n),\chi )`$, where $`\chi `$ is a primitive Dirichlet character of conductor $`2^n`$.

###### Theorem 4 (Emerton , Proposition 5.1)

Let $`m`$ be a positive integer greater than $`1`$, and let $`\chi `$ be a primitive Dirichlet character of conductor $`2^m`$. The smallest slope of the $`U_2`$ operator acting on cuspforms of weight $`k`$ and character $`\chi `$ is $`2^{3m}`$.

If we look at the character of conductor 4 and the odd character of conductor 8, there is a CM modular form which is defined over the field $`𝐐`$. We quote a result of Schoeneberg, proved in Ogg :

###### Theorem 5 (Ogg , Theorem VI.22)

Let $`i`$ be the square root of $`1`$, and let $`k`$ be an positive odd integer greater than $`3`$ and congruent to $`1mod4`$. Then there is a normalised cuspidal modular eigenform in $`S_k(\mathrm{\Gamma }_0(4),\tau )`$ with $`q`$-expansion

$$f_k(q)=\frac{1}{4}\underset{m,n𝐙}{}(m+ni)^{k1}q^{m^2+n^2}.$$

Let $`l`$ be a positive odd integer greater than $`1`$. Then there is a normalised cuspidal modular eigenform in $`S_l(\mathrm{\Gamma }_0(8),\tau \chi )`$ with $`q`$-expansion

$$g_l(q)=\frac{1}{2}\underset{m,n𝐙}{}(m+2ni)^{l1}q^{m^2+2n^2}$$

We see by inspection that the slope of $`f_k`$ is $`(k1)/2`$, and that the slope of $`g_l`$ is $`l1`$.

Hence we can, in certain cases, determine the smallest slope and another classical slope of $`U_2`$ acting on modular newforms of level $`\mathrm{\Gamma }_1(4)`$ or $`\mathrm{\Gamma }_1(8)`$ using previously known results.

Lawren Smithline has also proved results about the slopes of classical modular forms, and the techniques used in  are similar to those in this paper.

###### Theorem 6 (Smithline , Corollary 6.1.3.3)

Let $`v`$ be a non-negative integer and let $`k=23^{v+1}`$. Then there are exactly $`3^v`$ classical modular eigenforms of weight $`k`$ and level $`3`$ with $`3`$-slope $`3^{v+1}1`$.

Buzzard and Calegari  have proved the following theorem on the slopes of modular forms:

###### Theorem 7 (Buzzard-Calegari , Corollary 1)

The slopes of the $`U_2`$ operator on the space of cusp forms of weight $`0`$ are given by

$$1+2v_2\left(\frac{(3n)!}{n!}\right).$$

Jacobs  has proved, using similar techniques to those in this paper, that the slopes of the operator $`U_3`$ acting on the space of automorphic forms over a definite quarternion algebra are in the arithmetic progression $`\{1/2,3/2,5/2,\mathrm{}\}`$.

## 4 Overconvergent modular forms

A famous quote of Jacques Hadamard  says that “the shortest and best way between two truths of the real domain often passes through the imaginary one.” It seems that often the best way to prove results like Theorem 2 about classical modular forms is to prove a theorem for the overconvergent modular forms and then derive the theorem for classical modular forms as a consequence.

We therefore recall the definition of the 2-adic overconvergent modular forms, first by defining overconvergent modular forms of weight $`0`$, and then by deriving the definition for forms with weight and character.

Following Katz , section 2.1, we recall that, for $`E`$ be an elliptic curve over an $`𝐅_2`$-algebra $`R`$, there is a mod $`2`$ modular form $`A(E)`$ called the *Hasse invariant*, which has the $`q`$-expansion over $`𝐅_2`$ equal to $`1`$.

We consider the Eisenstein series of weight $`4`$ and tame level 1 defined over $`𝐙`$, with $`q`$-expansion

$$E_4(q):=1+240\underset{n=1}{\overset{\mathrm{}}{}}\left(\underset{0<d|n}{}d^3\right)q^n.$$

We see that $`E_4`$ is a lifting of $`A(E)^4`$ to characteristic 0, as the reduction of $`E_4`$ to characteristic $`2`$ has the same $`q`$-expansion as $`A(E)^4`$, and therefore $`E_4mod2`$ and $`A(E)^4`$ are both modular forms of level 1 and weight 4 defined over $`𝐅_2`$, with the same $`q`$-expansion.

Now the value of $`E_4(E)`$ may not be well-defined, but it can be shown that the valuation $`v_2(E_4(E))`$ is well-defined. This will allow us to define the ordinary locus of $`X_0(2^m)`$ and certain neighbourhoods of it. We follow the work of Coleman , and first define structures on the modular curve $`X_1(2^m)`$.

###### Definition 8 (Coleman , page 448)

Consider $`X_1(2^m)_{/𝐐_2}`$ as a rigid analytic space, and let $`t`$ be a point of $`X_1(2^m)`$.

If $`t`$ is a point of $`X_1(2^m)`$ which corresponds to a cusp, then we define $`v(E_4(t))=0`$, following , section 4.

We define the ordinary locus of $`X_1(2^m)`$ to be the set of points $`t`$ of $`X_1(2^m)`$ such that $`v_2(E_4(t))=0`$, and define $`Z_1(2^m)`$ to be the rigid connected component of the ordinary locus in $`X_1(2^m)`$ which contains the cusp $`\mathrm{}`$. This is a rigid analytic space.

In , page 36, it is shown that $`Z_1(2^m)`$ is an affinoid subdomain of the rigid space $`X_1(2^m)_{/𝐐_2}`$.

We will perform calculations in later sections on the modular curve $`X_0(2^m)`$, which we will now define.

###### Definition 9

Consider $`X_1(2^m)`$ as a modular curve. We see that the group $`G:=(𝐙/2^m𝐙)^\times `$ acts upon the non-cuspidal points of $`X_1(2^m)`$, by the following action: if $`a(𝐙/2^m𝐙)^\times `$, then the action of $`a`$ sends the pair $`(E,P)`$ to $`(E,aP)`$. This action extends to the cuspidal points of $`X_1(2^m)`$, and it sends cusps to cusps.

We will define the modular curve $`X_0(2^m)_{/𝐐_2}`$ to be the quotient of $`X_1(2^m)`$ by $`(𝐙/2^m𝐙)^\times `$.

We note that the action of the group $`G`$ does not change the valuation of a given elliptic curve $`E`$. We define $`Z_0(2^m)`$ to be the rigid connected component of the ordinary locus in $`X_0(2^m)`$ which contains the cusp $`\mathrm{}`$. It is a rigid analytic space.

We will now define strict affinoid neighbourhoods of $`Z_0(2^m)`$.

###### Definition 10 (Coleman , Section B2)

We think of $`X_0(2^m)`$ as a rigid space over $`𝐐_2`$, and we let $`tX_0(2^m)(\overline{𝐐}_2)`$ be a point, corresponding either to an elliptic curve defined over a finite extension of $`𝐐_2`$, or to a cusp. Let $`w`$ be a rational number, such that $`0<w<\mathrm{min}(2^{2m}/3,1/4`$).

We define $`Z_0(2^m)(w)`$ to be the connected component of the affinoid

$$\{tX_0(2^m):v_2(E_4(t))4w\}$$

which contains the cusp $`\mathrm{}`$.

The condition involves $`4w`$ rather than $`w`$ because we are working with a lifting of the fourth power of the Hasse invariant. Also, note that as $`E_4`$ is a lifting of the mod $`2`$ modular form $`A^4`$, and that any another lifting of $`A^4`$ would be of the form $`E_4+2F`$, where $`F`$ is a modular form, then this valuation is well-defined if $`0v(E_4(t))<1`$. This corresponds to the condition $`0w<1/4`$ in Definition 10.

We can now define overconvergent modular forms of weight $`0`$.

###### Definition 11 (Coleman, , page 397)

Let $`w`$ be a rational number, such that $`0<w<\mathrm{min}(2^{2m}/3,1/4)`$. Let $`𝒪`$ be the structure sheaf of $`Z_0(2^m)(w)`$.

We call sections of $`𝒪`$ on $`Z_0(2^m)(w)`$ *$`w`$-overconvergent 2-adic modular forms of weight $`0`$ and level $`\mathrm{\Gamma }_0(2^m)`$*.

If a section $`f`$ of $`𝒪`$ is a $`w`$-overconvergent modular form, then we say that $`f`$ is an *overconvergent $`2`$-adic modular form*.

Let $`K`$ be a complete subfield of $`𝐂_2`$, and define $`Z_0(2^m)(w)_{/K}`$ to be the affinoid over $`K`$ induced from $`Z_0(2^m)(w)`$ by base change from $`𝐐_2`$. The space $`M_0(2^m,w;K):=𝒪(Z_0(2^m)(w)_{/K})`$ of $`w`$-overconvergent modular forms of weight $`0`$ and level $`\mathrm{\Gamma }_0(2^m)`$ is a $`K`$-Banach space.

We now use non-cuspidal modular forms of the desired weight and character to define overconvergent modular forms with non-zero weight.

###### Definition 12

Let $`w`$ be a real number such that $`0<w<\mathrm{min}(2^{2m}/3,1/4)`$. Let $`k`$ be an integer and let $`\chi `$ be a character such that $`\chi (1)=(1)^k`$, and let $`K`$ be a complete subfield of $`𝐂_2`$. Let $`E_{k,\chi }`$ be the Eisenstein series of weight $`k`$ and character $`\chi `$.

The space of overconvergent 2-adic modular forms of weight $`k`$ and character $`\chi `$ is given by

$$_{k,\chi }(2^m,w;K):=E_{k,\chi }^{}_0(2^m,w;K).$$

This is a Banach space over $`K`$.

There are continuous Hecke operators $`T_p`$ which act on $`_{k,\chi }(2^m,w;K)`$; the operator $`U:=U_2`$, which is defined on $`q`$-expansions as

$$U\left(\underset{n=0}{\overset{\mathrm{}}{}}a_nq^n\right)=\underset{n=0}{\overset{\mathrm{}}{}}a_{2n}q^n$$

is also *compact* and therefore has a spectral theory.

As a consequence of results of Coleman, we have the following theorem:

###### Theorem 13 (Coleman , Theorem B3.2)

Let $`w`$ be a real number such that $`0<w<\mathrm{min}(2^{2m}/3,1/4)`$, let $`k`$ be an integer and let $`\chi `$ be a character such that $`\chi (1)=(1)^k`$.

The characteristic polynomial of $`U_2`$ acting on overconvergent 2-adic modular forms of weight $`k`$ and character $`\chi `$ is independent of the choice of $`w`$.

This theorem allows us to choose a convenient value of $`w`$ and prove results for that $`w`$, and guarantees that these results will hold for any $`w`$.

The connected component in Definition 10 is hard to work with. We will therefore rewrite it in terms of modular functions of level greater than $`1`$, to prove the following theorem:

###### Theorem 14

The spaces of overconvergent modular forms of weight 0 and level $`N=4\text{ or }8`$ are Tate algebras in one variable over $`𝐐_2(2^{4/N})`$.

We have given a valuation on the points $`t`$ of the rigid space $`X_0(2^m)`$, based on the lifting of the Eisenstein series $`E_4`$. We recall that the modular $`j`$-invariant is defined to be $`j=E_4^3/\mathrm{\Delta }`$. Therefore, we see that, if the elliptic curve corresponding to $`t`$ has good reduction, then $`\mathrm{\Delta }(t)`$ has valuation 0, and therefore that

$$v(t)=\frac{1}{4}v(E_4(t))=\frac{1}{12}v((E_4)^3(t))=\frac{1}{12}v(j(t)).$$

From Lemma 2.3 of , we see that there is a modular function $`j_8`$ which is a uniformiser on $`X_0(8)`$. It has $`q`$-expansion at $`\mathrm{}`$

$$j_8=\frac{1}{q_{n=1}^{\mathrm{}}(1+q^n)^4(1+q^{2n})^2(1+q^{4n})^4}=\left(\frac{\mathrm{\Delta }(q)\mathrm{\Delta }(q^4)}{\mathrm{\Delta }(q^2)\mathrm{\Delta }(q^8)}\right)^{1/12}.$$

Also, $`j_8(\mathrm{})=\mathrm{}`$.

There is another modular function $`j_{16}`$ which is a uniformiser on $`X_0(16)`$, with $`q`$-expansion at $`\mathrm{}`$ given by

$$j_{16}=\frac{1}{q_{n=1}^{\mathrm{}}(1+q^n)^2(1+q^{2n})(1+q^{4n})(1+q^{8n})^2}=\left(\frac{\mathrm{\Delta }(q^{16})^2\mathrm{\Delta }(q^2)}{\mathrm{\Delta }(q^8)\mathrm{\Delta }(q)^2}\right)^{1/24}.$$

We see also that $`j_{16}(\mathrm{})=\mathrm{}`$.

By an explicit calculation of $`q`$-expansions, using the formulae in Chapter 2 of , we see that

$$j=\frac{(j_8^4+256j_8^3+5120j_8^2+32768j_8+65536)^3}{(j_8^2+16j_8+64)(j_8+4)}$$

and

$$\frac{1}{j_8}=\frac{1}{j_{16}}+\frac{2}{j_{16}^2}.$$

Because we know that $`j_8(\mathrm{})=\mathrm{}`$, the connected component of $`Z_0(8)`$ which contains $`\mathrm{}`$ is of the form $`v(j_8)<D`$ for some rational number $`D`$. We see that, if $`v(j_8)<2`$, then $`v(j_8)=v(j)`$. This means that we have shown that

$$Z_0(8)(w)=\{xX_0(8):v(j_8(x))12w\}\text{ for }0<w<1/4.$$

Similarly, we see that the connected component of $`Z_0(16)`$ which contains $`\mathrm{}`$ is of the form $`v(j_8)<D`$ for some rational number $`D`$. We see that, if $`v(j_{16})<1`$, then $`v(j_8)=v(j_{16})`$, and therefore $`v(j_{16})=v(j_8)`$. This means that we have shown that

$$Z_0(16)(w)=\{xX_0(8):v(j_{16}(x))12w\},\text{ for }0<w<1/6.$$

We now define another modular function on $`X_0(2^m)`$, in terms of Eisenstein series.

Let $`k`$ be an integer, and let $`\theta :(𝐙/n𝐙)^\times 𝐂^\times `$ be a primitive Dirichlet character, such that $`\theta (1)=(1)^k`$. Recall from Washington , page 30, that the *extended Bernoulli numbers* $`B_{k,\theta }`$ are defined by

$$\underset{a=1}{\overset{N}{}}\frac{\theta (a)t\mathrm{exp}(at)}{\mathrm{exp}(Nt)1}=\underset{k=0}{\overset{\mathrm{}}{}}B_{k,\theta }\frac{t^i}{i!}$$

We define the Eisenstein series $`E_{k,\theta }^{}`$ to be

$$E_{k,\theta }^{}:=\frac{B_{k,\theta }}{2k}+\underset{n=1}{\overset{\mathrm{}}{}}\left(\underset{(d,p)=1}{\underset{0<d|n}{}}\theta (d)d^{k1}\right)q^n,$$

where $`B_{k,\theta }`$ is the extended Bernoulli number attached to $`k`$ and $`\theta `$.

There is an operator $`V`$ on the space of modular forms $`S_k(\mathrm{\Gamma }_0(N))`$; its effect on $`q`$-expansions is to send $`q`$ to $`q^2`$. We define $`V_k^{}`$ to be $`V(E_k^{})`$.

We define modular functions on $`X_0(8)`$ and $`X_0(16)`$ by

$$z_4:=\frac{E_{1,\tau }^{}/V_{1,\tau }^{}1}{2}=\frac{2}{j_8+2}$$

and

$$z_8:=\frac{E_{1,\chi \tau }^{}/V_{1,\chi \tau }^{}1}{\sqrt{2}}=\frac{\sqrt{2}}{j_{16}+2},$$

where we choose and fix a square root of $`2`$ in $`𝐂_2`$. These identities can be verified by explicit calculation.

Let $`w`$ be a rational number such that $`0<w<1/6`$. Then using the formulae above, we see that

$$Z_0(8)(w)=\{xX_0(8):v(z_4(x))112w\}.$$

We now choose $`w=1/12`$, to obtain

$$Z_0(8)(1/12)=\{xX_0(8):v(z_4(x))0\}.$$

Now, the rigid functions on the closed disc over $`𝐐_2`$ with centre 0 and radius 1 are defined to be power series of the form

$$\underset{n𝐍}{}a_nz^n:a_n𝐐_2,a_n0.$$

Therefore, the $`1/12`$-overconvergent modular forms of level $`\mathrm{\Gamma }_0(4)`$ and weight 0 are

$$𝐐_2z_4,$$

which is what we wanted to show. We now follow the same procedure for $`X_0(16)`$. Let $`w`$ be a rational number such that $`0<w<1/12`$. Then using the formulae above, we see that

$$Z_0(16)(w)=\{xX_0(16):v(z_8(x))1/212w\}.$$

We now choose $`w=1/24`$, to obtain

$$Z_0(16)(1/24)=\{xX_0(16):v(z_8(x))0\}.$$

Now, the rigid functions on the closed disc over $`𝐐_2`$ with centre 0 and radius 1 are defined to be power series of the form

$$\underset{n𝐍}{}a_nz^n:a_n𝐐_2(\sqrt{2}),a_n0.$$

Therefore, the $`1/24`$-overconvergent modular forms of level $`\mathrm{\Gamma }_0(8)`$ and weight 0 are

$$𝐐_2(\sqrt{2})z_8,$$

so we have shown that these spaces of modular forms are Tate algebras in one variable. $`\mathrm{}`$

We now show that the odd powers of $`z_N`$ are sent to $`0`$ under the $`U_2`$ operator.

###### Lemma 15

Let $`N=4`$ or $`8`$ and let $`i`$ be a positive integer. Then

$$U_2(z_N^{2i+1})=0,\text{ and }U(z_N^{2i})=\left(U(z_N^2)\right)^i.$$

Recall that the Eisenstein series $`E_𝐤^{}`$ is an eigenvector with eigenvalue 1 for $`U_2`$, and that $`U_2(V(E_𝐤^{}))=E_𝐤^{}`$. Then we see that we have (for $`\mu =2\text{ or }\sqrt{2}`$):

$`U_2(z)`$ $`=`$ $`U_2\left({\displaystyle \frac{E_𝐤^{}/V_𝐤^{}1}{\mu }}\right)={\displaystyle \frac{1}{\mu }}U_2\left({\displaystyle \frac{E_𝐤^{}V_𝐤^{}}{V_𝐤^{}}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{\mu E_𝐤^{}}}U_2(E_𝐤^{}V_𝐤^{})={\displaystyle \frac{1}{\mu E_𝐤^{}}}(E_𝐤^{}E_𝐤^{})=0.`$
Hence we see that $`z_N`$ has only odd $`q`$-coefficients, and that therefore $`z_N^2`$ has only even $`q`$-coefficients.

Therefore $`z_N^{2i+1}`$ has only odd $`q`$-coefficients. Hence for all non-negative integers $`t`$, we see that

$$U_2\left(z_N^{2t+1}\right)=0.$$

Because we have just shown that $`z_N`$ has only odd $`q`$-coefficients, we see that

$$z_N=qF(q^2)=qV(F(q)),$$

for some power series $`F(q)`$.

Therefore we have

$$U_2(z_N^{2i})=U_2(q^{2i}V(F(q)^{2i}))=U_2(V(q^iF(q)^{2i})),$$

and hence we see that

$$U_2(z_N^{2i})=q^iF(q)^{2i}=\left(qF(q)^2\right)^i=U_2(z_N^2)^i,$$

which proves the Lemma. $`\mathrm{}`$

Because we have written down the overconvergent modular forms as an explicit Banach space, we can write down its *Banach basis*: the set $`\{z_4,z_4^2,z_4^3,\mathrm{}\}`$ forms a Banach basis for the overconvergent modular forms of weight $`0`$ and level $`\mathrm{\Gamma }_0(4)`$ and the functions $`\{z_8,z_8^2,z_8^3,\mathrm{}\}`$ form a Banach basis for the overconvergent modular forms of weight $`0`$ and level $`\mathrm{\Gamma }_0(8)`$.

These Banach bases are composed of weight $`0`$ modular functions — we want to be able to consider the action of the $`U`$ operator on overconvergent modular forms with non-zero weight-character $`𝐤`$. Using an observation from the work of Coleman, we will be able to move between weight-character $`\mathrm{𝟎}`$ and weight-character $`𝐤`$ via multiplication by a suitable quotient of modular forms.

From the discussion in Coleman , page 450, we see that the $`U_2`$ operator acting on overconvergent modular forms of weight-character $`(kt,\theta ^t)`$ where $`t`$ is odd has the same characteristic power series as the composition of the $`U`$ operator acting on overconvergent modular forms of weight-character $`\mathrm{𝟎}`$ with multiplication by the $`z`$-expansion of $`(E_{k,\theta }^{}/V_{k,\theta }^{})^t`$.

For level $`4`$ and weight $`2t+1`$ and level $`8`$ and weight either of the form $`2t+1`$ or $`4k+2`$, we can write the weight $`k`$ multiplier as a power of $`1+2z_4`$ or $`1+\sqrt{2}z_8`$. If $`N=8`$ and the weight $`k`$ divides $`4`$, then we can write $`E_{k,\theta }^{}/V_{k,\theta }^{}`$ as a rational function of $`z_8`$, because $`z_8`$ is a rational function of the uniformiser $`j_{16}`$ of the genus $`0`$ modular curve $`X_0(16)`$. For instance, we see that

$$\frac{E_{4,\chi }^{}}{V_{4,\chi }^{}}=\frac{11+2z_8+24z_8^248z_8^316z_8^4352z_8^5}{11+24z_8^216z_8^4}.$$

We can consider the action of the $`U_2`$ operator on these spaces of overconvergent modular forms.

###### Definition 16

Let $`M=(m_{i,j})`$ be the infinite compact matrix of the operator $`U_2(E_𝐤^{}/V_𝐤^{})^t`$ acting on the Banach basis $`\{z_N,z_N^2,\mathrm{}\}`$, where $`m_{i,j}`$ is defined to be the coefficient of $`z_N^j`$ in the $`z_N`$-expansion of $`U_2(z_N^i)(E_𝐤^{}/V_𝐤^{})^t`$.

Because we know that $`_0`$ is a Banach space, we can show that the matrix $`M`$ is a compact matrix; in other words, the trace, determinant and characteristic power series of $`M`$ are all well-defined.

We will use a theorem of Serre to prove our theorem on the slopes of $`U`$ acting on $`_𝐤`$.

###### Theorem 17 (Serre , Proposition 7)

1. Let $`M`$ be an $`n\times n`$ matrix defined over a finite extension of $`𝐐_2`$, and let $`0r𝐐`$. Let $`det(1tM)=_{i=0}^nc_it^i`$. Let $`M_m`$ be the matrix formed by the first $`m`$ rows and columns of $`M`$.

Assume that

1. For all positive integers $`m`$ such that $`1mn`$, the valuation of $`det(M_m)`$ is $`r\frac{m(m+1)}{2}`$.
    2. The valuation of elements in column $`j`$ is at least $`rj`$.

Then we have that, for all positive integers $`m`$ such that $`1mn`$, $`v_2(c_m)=r\frac{m(m+1)}{2}`$.
2. Let $`M_{\mathrm{}}`$ be a compact infinite matrix. If there is a sequence of finite matrices $`M_m`$ which tend to $`M_{\mathrm{}}`$, then the finite characteristic power series $`det(1tM_m)`$ converge coefficientwise to $`det(1tM_{\mathrm{}})`$, as $`m\mathrm{}`$.

We now quote a result of Coleman that tells us that overconvergent modular forms of small slope are in fact classical modular forms:

###### Theorem 18 (Coleman , Theorem 1.1)

Let $`k`$ be a non-negative integer. Every $`2`$-adic overconvergent modular eigenform of weight $`k`$ with slope strictly less than $`k1`$ is a classical modular form.

We will now state the theorem on the slopes of overconvergent modular forms of weight-character $`(2k+1,\tau )`$ and weight-character $`(k,\chi \tau ^k)`$. This, combined with Theorem 18 and a knowledge of the dimensions of spaces of classical cusp forms, will suffice to prove Theorem 2 and Corollary 3.

###### Theorem 19

Let $`k`$ be an integer, let $`\tau `$ be the primitive Dirichlet character of conductor $`4`$ and let $`\chi `$ be the even primitive Dirichlet character of conductor $`8`$.

1. The slopes of overconvergent 2-adic cuspidal eigenforms of weight $`2k+1`$ and character $`\tau `$ are $`\left\{2i\right\}_{i𝐍}`$.
2. The slopes of overconvergent 2-adic cuspidal eigenforms of weight $`k`$ and character $`\chi \tau ^k`$ are $`\left\{i\right\}_{i𝐍}`$.

We now recall a theorem of Cohen and Oesterlé:

###### Theorem 20 (Cohen-Oesterlé , Théorème 1)

Let $`\chi `$ be a primitive Dirichlet character of conductor $`2^m>2`$, and let $`k`$ be a positive integer. Let $`𝐤=(k,\chi )`$ be an integral weight-character.

The dimension of the space of cuspidal modular forms of weight-character $`𝐤`$ is

$$2^{m3}(k1)1.$$

We see that the slopes of the first $`2^{m3}(k1)1`$ overconvergent modular forms of level $`\mathrm{\Gamma }_0(2^m)`$ and primitive Dirichlet character are

$$2^{3m},2^{2(3m)},\mathrm{},k12^{3m}.$$

Therefore, using Theorem 18, we see that all of these slopes are classical, because $`k12^{3m}<k1`$. Hence we have proved Corollary 3. $`\mathrm{}`$

## 5 Necessary conditions on matrices

We will show that we can apply Theorem 17 by proving the following theorem:

###### Theorem 21

Let $`N=4\text{ or }8`$, and let $`M`$ be the matrix of the $`U`$-operator acting on the Banach basis $`\left\{(z_N)^i\right\}`$. Define the set $`a_i`$ by $`U(z_N^2):=_{i=1}^{\mathrm{}}a_i(z_N)^i`$. Assume that

$$\mathrm{}:v(a_i)=i4/N,\mathrm{for}i\mathrm{odd},\mathrm{and}v(a_i)>i4/N,\mathrm{for}i\mathrm{even}.$$

Then the valuation of the determinant of the $`i\times i`$ matrix $`M_n`$ is $`8i/N`$.

We can show the precondition of this theorem directly, by considering the two identities of modular functions

$$U(z_4^2)=\frac{2z_4}{(1+2z_4)^2}\text{ and }U(z_8^2)=\frac{\sqrt{2}z_8}{1+2z_8^2},$$

where we have chosen a square root of $`2`$ in the extension $`𝐐_2(\sqrt{2})`$. These can be verified by multiplying out both sides and transforming them into an identity of modular *forms*; we then use a theorem of Sturm  to show that both sides of the equation are the same by checking the $`q`$-expansions at $`\mathrm{}`$.

We then view these rational functions of $`z_N`$ as generating functions for power series in $`z_N`$.

We will also prove a result about the ring over which the $`q`$-expansions at $`\mathrm{}`$ of the cusp forms of weight-character $`𝐤`$ are defined. This will allow us to prove Corollary 3.

###### Corollary 22

Let $`𝐤=(k,\chi )`$ be an integral weight-character, where $`\chi `$ is a primitive Dirichlet character of conductor $`4`$ or $`8`$ and $`k`$ is a positive integer such that $`\chi (1)=(1)^k`$. Let $`K`$ be the field $`𝐐_2`$ if $`N=4`$ or $`𝐐_2(\sqrt{2})`$ if $`N=8`$.

Let $`f=_{n=1}^{\mathrm{}}a_nq^n`$ be a normalised classical cuspidal modular eigenform of weight-character $`𝐤`$.

Then $`a_nK`$ for all $`n`$.

Theorem 19 tells us that the slopes of overconvergent modular forms of weight-character $`𝐤`$ are distinct. Hence by applying Theorem 20 and Theorem 18 we see that the slopes of classical modular forms of weight-character $`𝐤`$ are also distinct.

We will now recall a fact from Ribet , page 21, to prove this Corollary. Let $`\sigma `$ be an element of $`Gal(\overline{K}/K)`$. Then we have that

$$\sigma (f):=\underset{n=1}{\overset{\mathrm{}}{}}\sigma (a_n)q^n$$

is a classical cuspidal modular eigenform of weight-character $`𝐤`$.

We see that the valuation of $`\sigma (a_2)`$ is the same as that of $`a_2`$, because the characteristic polynomial of $`a_2`$ is stable under conjugation by $`\sigma `$. Therefore, $`\sigma (f)`$ is an eigenform of weight-character $`𝐤`$ with the same slope as $`f`$. Hence $`\sigma (f)=f`$, because there is only one classical eigenform of weight-character $`𝐤`$ which has any given slope, by Theorem 18 and Theorem 19.

This means that $`\sigma (f)=f`$ for all $`\sigma `$. Therefore $`a_nK`$ for all positive integers $`n`$. $`\mathrm{}`$

We will perform a series of transformations to obtain a matrix which we can apply Theorem 17 to.

We first note that the odd-numbered columns of the matrix $`M_n`$ are identically zero, because we have shown that $`U(z_N^{2a+1})=0`$.

We consider the matrix $`O_n`$, defined by

$$\left(O_n\right)_{i,j}:=\left(M_{2n}^{}\right)_{2i,2j},\mathrm{where}\mathrm{\hspace{0.33em}1}i,jn.$$

This has the same characteristic power series as $`M_{2n}`$.

We now show that the matrices $`O_n`$ have determinant of valuation $`8n/N`$, in order to be able to use Theorem 17.

To do this we will pre- and post-multiply the matrix $`O_n`$ by diagonal matrices to obtain a matrix $`O^{}`$ which has elements of valuation at least $`8i/N`$ in column $`i`$. Let $`D(\alpha )`$ be the diagonal matrix with $`(i,i)^{th}`$ coefficient $`\alpha ^i`$. We let $`\alpha `$ be $`2`$ if $`N=4`$ and a square root of $`2`$ if $`N=8`$, and we define

$$O^{}=D(\alpha ^1)O_nD(\alpha ).$$

It can now be checked that the valuation of the elements in the $`i^{th}`$ column of $`O^{}`$ is at least $`8i/N`$.

We will show that the matrix $`O^{}`$ has determinant with valuation $`8/Nn(n+1)`$ by showing that it is the product of two matrices, one of which is the diagonal matrix $`D(\alpha )`$, and one of which is a matrix with determinant a unit.

We define the matrix $`P`$ to be $`D(\alpha )^1O^{}`$. The entries of $`P`$ are given by $`P_{i,j}=\alpha ^iO_{i,j}^{}`$ and are therefore elements of the ring of integers of $`𝐐_2(\sqrt{2})`$, because the valuation of elements in the $`i^{th}`$ column of $`O^{}`$ is at least $`i`$. Therefore, we can define the matrix $`P^{}`$ by reducing the entries of $`P`$ modulo $`2`$; if $`P^{}`$ has determinant $`1`$ in $`𝐅_2`$ then $`O^{}`$ has determinant $`8/Nn(n+1)`$.

By tracing the definition of the columns through all of this, we see that the odd-indexed columns of the mod $`2`$ matrix in weight $`(2i+1)k`$ have generating functions

$$\mathrm{col}(2l+1)=\frac{x^{l+1}y^{2l+1}(1+x)^i}{(1+x)^{2l+1}},$$

and the even-indexed columns of the matrix have generating functions

$$\mathrm{col}(2l)=\frac{x^ly^{2l}(1+x)^i}{(1+x)^{2l}}.$$

This is clear from the discussion earlier if $`N=4`$ or $`N=8`$ and $`4`$ does not divide the weight.

If $`4`$ does divide the weight, then we consider the quotient $`E_{k,\chi }^{}/V_{k,\chi }^{}`$; notice that the $`q`$-expansions of the Eisenstein series $`E_{2,\chi }^{}`$ and $`E_{k,\chi }^{}`$ are congruent modulo $`2`$, by Fermat’s little theorem. Let $`z_k:=(E_{k,\chi }^{}/V_{k,\chi }^{}1)/\lambda _k`$, where $`\lambda _k`$ is a square root of $`2k/B_{k,\chi }`$. This is also an overconvergent modular form, so it is an element of $`𝐐_2(\sqrt{2})z_8`$. Then the effect of the changes above on the columns is to send the multiplier $`(1+\lambda _kz_k)`$ to $`1+z_k`$. We see that $`z_kz_8mod8`$ because

$$z_k=\frac{E_{k,\chi }^{}V_{k,\chi }^{}}{\lambda _kV_{k,\chi }^{}}E_{k,\chi }^{}V_{k,\chi }^{}E_{2,\chi }^{}V_{2,\chi }^{}\frac{E_{2,\chi }^{}V_{2,\chi }^{}}{\sqrt{2}V_{2,\chi }^{}}z_8.$$

Therefore the multiplier $`E_{k,\chi }^{}/V_{k,\chi }^{}`$ when expanded as a rational function of $`z_8`$ and then altered by the procedure above is congruent modulo $`2`$ to $`1+z_8`$.

For example, we consider the $`z_8`$ expansion of $`E_{4,\chi }^{}/V_{4,\chi }^{}`$. After sending $`2z_8`$ to $`z_8`$, we see that we have

$$\frac{11+z_8+6z_8^26z_8^3z8^411z_8^5}{11+6z_8^2z_8^4}\frac{1+z_8+z8^4+z8^5}{1+z_8^4}1+z_8mod2.$$

I would like to thank Robin Chapman  for the idea behind the following proof. To show that the $`n\times n`$ mod $`2`$ matrix has determinant $`1`$, we will show that the elements $`C:=\left\{\mathrm{col}^{}(1)=\mathrm{col}1(1+x)^N,\mathrm{},\mathrm{col}^{}(N)=\mathrm{col}N(1+x)^N\right\}`$ are a basis of the ring $`𝐅_2[x]/(x^{N+1})`$. Because $`(1+x)`$ is a unit in the ring, we may multiply each $`\mathrm{col}(i)`$ by $`(1+x)^{Ni}`$ to make the calculations easier. We see that there are exactly the right number of elements, so we must check that they are linearly independent.

We write $`_{i=1}^N\lambda _i\mathrm{col}^{}(i)=0`$. We will show that all of the $`\lambda _i=0`$, so that the columns are linearly independent.

We see that the element $`\mathrm{col}^{}(1)=x(1+x)^{N1}`$ is the only element of the set $`C`$ which has an $`x^N`$ term. Therefore $`\lambda _1=0`$. Also, we see that $`\mathrm{col}^{}(2)=x(1+x)^{N2}`$ is now the only nonzero term with an $`x`$ term, so therefore $`\lambda _2=0`$.

By continuing this process, we show that each $`\lambda _i`$ must be zero. Therefore the set $`C`$ is composed of linearly independent elements and hence it is a basis. So the determinant of the mod $`2`$ matrix is $`1`$. Hence we have shown Theorem 21 and therefore Theorem 19. $`\mathrm{}`$

## 6 Acknowledgements

I would like to thank Kevin Buzzard for many helpful conversations and much inspiration while he was my PhD advisor. I would also like to thank Edray Goins and Ken McMurdy for helpful conversations, and Robin Chapman for his suggestion on a shorter proof of Theorem 21. I would also like to thank David Burns and Frazer Jarvis for their careful reading of my PhD thesis , which is where the results of this paper first appeared.

Some of the computer calculations were executed by the computer algebra package magma  running on the machine crackpipe, which was bought by Kevin Buzzard with a grant from the Central Research Fund of the University of London. I would like to thank the CRF for their support. Some other calculations were performed on the machine shimura, which is based at the University of California at Berkeley. I would like to thank the administrators, Wayne Whitney and William Stein, for allowing me to hold an account on their machine.