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13159812	https://en.wikipedia.org/wiki/Borassus%20flabellifer	Borassus flabellifer	Borassus flabellifer, commonly known as doub palm, palmyra palm, tala or tal palm, toddy palm, lontar palm, wine palm, or ice apple, is a fan palm native to South Asia (especially in Bangladesh, East India, and South India) and Southeast Asia. It is reportedly naturalized in Socotra. Description Borassus flabellifer is a robust tree and can reach a height of . The trunk is grey, robust, and ringed with leaf scars; old leaves remain attached to the trunk for several years before falling cleanly. The leaves are fan-shaped, and long, with robust black teeth on the petiole margins. Like all Borassus species, B. flabellifer is dioecious with male and female flowers on separate plants. But very rarely male and female flowers in same trees have also been noticed The male flowers are less than long and form semi-circular clusters, which are hidden beneath scale-like bracts within the catkin-like inflorescences. In contrast, the female flowers are golfball-sized and solitary, sitting upon the surface of the inflorescence axis. After pollination, these blooms develop into fleshy fruits wide, each containing 1-3 seeds. The fruits are black to brown with sweet, fibrous pulp, and each seed is enclosed within a woody endocarp. Young palmyra seedlings grow slowly, producing only a few leaves each year (establishment phase), but at an as yet undetermined time, they grow rapidly, producing a substantial stem. Uses Fruit The fruit (palmyra fruit) measures to in diameter, has a black husk, and is borne in clusters. The top portion of the fruit must be cut off to reveal the sweet jelly seed sockets, translucent pale-white, similar to that of the lychee but with a milder flavor and no pit. The sweet jelly seed sockets occur in combinations of two, three or four seeds inside the fruit. The jelly part of the fruit is covered with a thin, yellowish-brown skin. These are known to contain watery fluid inside the fleshy white body. These seed sockets have been the inspiration behind certain sandeshes called jolbhora (জলভরা) found in Bengal. The soft orange-yellow mesocarp pulp of the ripe fruit is sugary, dense and edible, rich in vitamins A and C. They also contain bitter compound called flabelliferrins, which are steroidal saponins. The conventional way this fruit is eaten is when the outer casing is still unripe while the seeds are eaten as the fruit. But if the entire fruit is left to ripen, the fibrous outer layer of the palm fruits can also be eaten raw, boiled, or roasted. When this happens, the fruit takes a purple-blackish hue, and tastes similar to coconut flesh. The skin is also eaten as part of the fruit similar to how mango skins are often consumed along with the fruit. Bengalis have perfected the art of making various sweet dishes with the yellowish viscous fluid substance obtained from a ripe palm fruit. These include mustard oil-fried (alternately sunflower oil-fried) taler bora (তালের বড়া) "palmyra vadas" or mixed with thickened milk to prepare tal-khir (তাল ক্ষীর). Thais also use the fruit to make the steamed fluffy tala palm cake, call “Khanom Tan”. In northern India, the fruit is known as Taad Gola in Hindi-Urdu (ताड़ गोला / ). In Kerala it is called nonku (നൊങ്ക്) whereas in Tamil Nadu, it is called nungu (நுங்கு). In Odisha, it is called tala (ତାଳ). Ice apple in Indonesia is called buah lontar or siwalan. In Karnataka it is called "Taati Nungu"(ತಾಟಿ ನುಂಗು / ತಾಟಿ ನಿಂಗು). In Myanmar, it is called htan-thee (ထန်းသီး). In Telangana and Andhra Pradesh, this fruit is called as "Thaati Munjalu" (తాటి ముంజలు). In Tulu language of Coastal Karnataka it is called “Erolu”(ಇರೋಲು). Sap Obtaining the sap traditionally involves tapping the top shoots and collecting the dripping juice in hanging earthen pots (in some regions a plastic or bamboo bottle). The juice collected in evening or after fermentation becomes sour, and is called tāḍī (ताडी > "toddy") in Marathi, hta-yay (ထန်းရည်) in Myanmar and Bhojpuri. This sap was the main source of sugar production in Thailand before sugarcane was introduced, as can be seen in the Thai word for sugar (), which literally means the water of the tala palm. A sugary sap can be obtained from the young inflorescence, either male or female and it is concentrated to a crude sugar called jaggery or Tal Patali (তাল পাটালী) in Bengali, hta-nyat (ထန်းလျက်) in Myanmar and Pana Vellam or Karuppukatti (கருப்புகட்டி or கருபட்டி) in Tamil, or it can be fermented to make an alcoholic beverage called toddy or htan-yay hkar (ထန်းရည်ခါး) in Myanmar, or distilled to make a liqour arrack. The concentrated raw sugar obtained from palms is called Gula Jawa (Javanese sugar) in Indonesia, and is widely used in Javanese cuisine. In Thailand, it is called nam tan pik (น้ำตาลปึก), referring to the pack of sugar obtained from drying the palm sap, though in the modern day nam tan pik is often made from coconut water because the convenient of farming and harvesting. In Thailand, there are techniques that utilize the anti-bacterial agents of some woods to keep the sap from becoming sour while tapping. After sterilization, the sap is available as a beverage called nam tan sod (น้ำตาลสด, ) or used to make an alcoholic beverage called nam tan mao (น้ำตาลเมา ). “Nam tan”, literally means tala palm water, later acquired the meaning of “sugar “. Sprouts In the Indian states of Tamil Nadu, Andhra Pradesh, Telangana and Bihar, and in Jaffna, Bengal, Sri Lanka, the seeds are planted and made to germinate and the fleshy stems (below the surface) are boiled or roasted and eaten. It is very fibrous and nutritious. It is known as Thegalu (తేగలు) or Gaygulu (గేగులు) or Gengulu (గెంగులు) (especially in Telangana) in Telugu, as Panai Kizhangu or Panangkizhangu (பனங்கிழங்கு) in Tamil, and as htabin myiq (ထန်းပင်မြစ်) in Myanmar. The germinated seed's hard shell is also cut open to take out the crunchy kernel, which tastes like a sweeter water chestnut. It is called "Taal-Anti" (তাল আঁটি) in Bengali, "Sachi-Htway" (ဆံချည်ထွေး) in Myanmar, "Buragunju" (బురగుంజు) in Telugu and "Thava nai" in Tamil. The white kernel of the ripe palm fruit after being left for a few months is used as an offering in Lakshmi Puja in various parts of Bengal and is also eaten raw. In Thai cuisine, it is used as an ingredient to a type of curry, called “Kaeng Hua Tan”. Leaves The Borassus flabellifer leaves are used for thatching, mats, baskets, fans, hats, umbrellas, and as writing material. All the literature of the old Tamil was written in preserved palm leaves also known as Palm-leaf manuscript. In Tamil Yaedu or Olai chuvadi. Most of the ancient literature in Telugu are written on palm leaves (Tala patra grandhas). In Indonesia the leaves were used in the ancient culture as paper, known as "lontar" (from Old/Modern Javanese ron tal "tal leaves") Leaves of suitable size, shape, texture, and maturity were chosen and then seasoned by boiling in salt water with turmeric powder, as a preservative. The leaves were then dried. When they were dry enough, the face of the leaf was polished with pumice, cut into the proper size, and a hole made in one corner. Each leaf made four pages. The writing was done with a stylus and had a very cursive and interconnected style. The tal is so closely related to regional manuscript culture that a tal frond is immortalized as the part of the logo for the Dewan Bahasa dan Pustaka, Malaysia's language regulatory board designed in 1957 by Hussien Enas. The stem of the leaves has thorny edges (called "karukku" in Tamil). The skin of the stem can be peeled off and be used as rope and also used to weave into cots (நார்க்கட்டில் in Tamil). In some part of Tamil Nadu, a variety of rice flour cake (called "Kozhukattai") is prepared using the leaf. In the eastern part of India, the leaves are used to make hand fans. In Myanmar, the leaves are used to make hand fans for the Buddhist monks and are called "Yap" (ယပ်). Trunk The stalks are used to make fences and also produce a strong, wiry fiber suitable for cordage and brushes. The black timber is hard, heavy, and durable and is highly valued for construction. It is superior to coconut timber, or red palm. Crown When the crown of the tree is removed, the segment from which the leaves grow out is an edible cake. This is called htan-ohn-hnauk (ထန်းဦးဏှောက်) in Myanmar, pananchoru (பனஞ்சோறு) in Tamil or thati adda (తాటి అడ్డ/తాటి మట్ట) in Telugu. Roots In Cambodia, where the palm is known as thnôt''' (Khmer), the roots are dried and smoked to heal nasal complaints. CultivationBorassus flabellifer has a growth pattern, very large size, and clean habits that make it an attractive ornamental tree, cultivated for planting in gardens and parks as landscape palm species. Cultural symbolism The palmyra tree is the official tree of Tamil Nadu. Highly respected in Tamil culture, it is called "katpaha tharu" ("celestial tree") because all its parts have a use. Panaiveriyamman, named after panai, the Tamil name for the Palmyra palm, is an ancient tree deity related to fertility linked to this palm. This deity is also known as Taalavaasini, a name that further relates her to all types of palms. The Asian palmyra palm is a symbol of Cambodia where it is a very common palm, found all over the country. It also grows near the Angkor Wat temple. In Indonesia the Palmyra tree is the symbol of South Sulawesi province. This plant has captured the imagination of Bengalis, especially in the words of Rabindranth Tagore whose nursery rhyme 'Tal Gach ek Paye dariye' (তাল গাছ এক পায়ে দাড়িয়ে.., literally Palmyra tree standing on a single leg ... ) in Sahaj Path (সহজ পাঠ) is a staple reading material in schools in Bangladesh and West Bengal. In the Hindu epic Mahabharata'', a palmyra tree is the chariot-banner of Bheeshma and Balarama. In Myanmar, the tree is the symbol of Anyar (အညာ) (the dry zone of Myanmar), and is called "pa-de-thar-pin" (ပဒေသာပင်) meaning the tree from which anything you wish can be taken. There are many poems and traditional sounds related to this tree. Sunthorn Phu, Thailand’s eminent bard of the Early Bangkok Era, mentioned the plant in many of his poems.	Biology and health sciences	Arecales (inc. Palms)	Plants
13163733	https://en.wikipedia.org/wiki/Stenotherm	Stenotherm	A stenotherm (from Greek στενός stenos "narrow" and θέρμη therme "heat") is a species or living organism capable of surviving only within a narrow temperature range. This specialization is often found in organisms that inhabit environments with relatively stable environments, such as deep sea environments or polar regions. The opposite of a stenotherm is a eurytherm, an organism that can function across a wide range of body temperatures. Eurythermic organisms are typically found in environments with significant temperature variations, such as temperate or tropical regions. The size, shape, and composition of an organism's body can influence its temperature regulation, with larger organisms generally maintaining a more stable internal temperature than smaller ones. Examples Chionoecetes opilio is a stenothermic organism, and temperature significantly affects its biology throughout its life history, from embryo to adult. Small changes in temperature (< 2 °C) can increase the duration of egg incubation for C. opilio by a full year.	Biology and health sciences	Basics	Biology
14329297	https://en.wikipedia.org/wiki/Climate%20change%20in%20Australia	Climate change in Australia	Climate change has been a critical issue in Australia since the beginning of the 21st century. Australia is becoming hotter and more prone to extreme heat, bushfires, droughts, floods, and longer fire seasons because of climate change. Climate issues include wildfires, heatwaves, cyclones, rising sea levels, and erosion. Since the beginning of the 20th century, Australia has experienced an increase of over 1.5 °C in average annual temperatures, with warming occurring at twice the rate over the past 50 years compared with the previous 50 years. Recent climate events such as extremely high temperatures and widespread drought have focused government and public attention on the effects of climate change in Australia. Rainfall in southwestern Australia has decreased by 10–20% since the 1970s, while southeastern Australia has also experienced a moderate decline since the 1990s. Rainfall is expected to become heavier and more infrequent, as well as more common in summer rather than in winter. Australia's annual average temperatures are projected to increase 0.4–2.0 °C above 1990 levels by the year 2030, and 1–6 °C by 2070. Average precipitation in the southwest and southeast Australia is projected to decline during this time, while regions such as the northwest may experience increases in rainfall. Climate change is affecting the continent's environment and ecosystems. Australia is vulnerable to the effects of global warming projected for the next 50 to 100 years because of its extensive arid and semi-arid areas, and already warm climate, high annual rainfall variability. The continent's high fire risk increases this susceptibility to changes in temperature and climate. Meanwhile, Australia's coastlines will experience erosion and inundation from an estimated increase in global sea level. Australia's unique ecosystems such as the Great Barrier Reef and many animal species are also at risk. Climate change also has diverse implications for Australia's economy, it's agriculture and public health. Projected impacts include more severe floods, droughts, and cyclones. Furthermore, Australia's population is highly concentrated in coastal areas at risk from rising sea levels, and existing pressures on water supply will be exacerbated. The exposure of Indigenous Australians to climate change impacts is exacerbated by existing socio-economic disadvantages which are linked to colonial and post-colonial marginalisation. The communities most affected by climate changes are those in the North where Aboriginal and Torres Strait Islander people make up 30% of the population. Aboriginal and Torres Strait Islander communities located in the coastal north are the most disadvantaged due to social and economic issues and their reliance on traditional land for food, culture, and health. This has raised the question for many community members in these areas, "Should we stay or move away?" Australia is also a contributor to climate change, with its greenhouse gas emissions per capita above the world average. The country is highly reliant on coal and other fossil fuels, although renewable energy coverage is increasing. National climate change mitigation efforts include a commitment to achieving net zero emissions by 2050 under the Paris Agreement, although Australia has repeatedly ranked poorly in the Climate Change Performance Index and other international rankings for its climate targets and implementation. Climate change adaptation can be performed at national and local levels and was identified as a priority for Australia in the 2007 Garnaut Review. Climate change has been a divisive or politicised issue in Australian politics since the 2000s, contributing to successive governments implementing and repealing mitigation policies such as carbon pricing. Some Australian media outlets have promoted climate misinformation. The issue has sparked protests in support of climate change policies, including some of the largest demonstrations in Australia's history. Greenhouse gas emissions Impacts on the natural environment Temperature and weather changes Australia's instrumental record from 1885 to the present shows the following broad picture: Conditions from 1885 to 1898 were generally fairly wet, though less so than in the period since 1968. The only noticeably dry years in this era were 1888 and 1897. Although some coral core data suggest that 1887 and 1890 were, with 1974, the wettest years across the continent since settlement, rainfall data for Alice Springs, then the only major station covering the interior of the Northern Territory and Western Australia, strongly suggest that 1887 and 1890 were overall not as wet as 1974 or even 2000. In New South Wales and Queensland, however, the years 1886–1887 and 1889–1894 were indeed exceptionally wet. The heavy rainfall over this period has been linked with a major expansion of the sheep population and February 1893 saw the disastrous 1893 Brisbane flood. A drying of the climate took place from 1899 to 1921, though with some interruptions from wet El Niño years, especially between 1915 and early 1918 and in 1920–1921, when the wheat belt of the southern interior was drenched by its heaviest winter rains on record. Two major El Niño events in 1902 and 1905 produced the two driest years across the whole continent, whilst 1919 was similarly dry in the eastern States apart from the Gippsland. The period from 1922 to 1938 was exceptionally dry, with only 1930 having Australia-wide rainfall above the long-term mean and the Australia-wide average rainfall for these seventeen years being 15 to 20 per cent below that for other periods since 1885. This dry period is attributed in some sources to a weakening of the Southern Oscillation and in others to reduced sea surface temperatures. Temperatures in these three periods were generally cooler than they are currently, with 1925 having the coolest minima of any year since 1910. However, the dry years of the 1920s and 1930s were also often quite warm, with 1928 and 1938 having particularly high maxima. The period from 1939 to 1967 began with an increase in rainfall: 1939, 1941 and 1942 were the first close-together group of relatively wet years since 1921. From 1943 to 1946, generally dry conditions returned, and the two decades from 1947 saw fluctuating rainfall. 1950, 1955 and 1956 were exceptionally wet except 1950 and 1956 over arid and wheatbelt regions of Western Australia. 1950 saw extraordinary rains in central New South Wales and most of Queensland: Dubbo's 1950 rainfall of can be estimated to have a return period of between 350 and 400 years, whilst Lake Eyre filled for the first time in thirty years. In contrast, 1951, 1961 and 1965 were very dry, with complete monsoon failure in 1951/1952 and extreme drought in the interior during 1961 and 1965. Temperatures over this period initially fell to their lowest levels of the 20th century, with 1949 and 1956 being particularly cool, but then began a rising trend that has continued with few interruptions to the present. Since 1968, Australia's rainfall has been 15 per cent higher than between 1885 and 1967. The wettest periods have been from 1973 to 1975 and 1998 to 2001, which comprise seven of the thirteen wettest years over the continent since 1885. Overnight minimum temperatures, especially in winter, have been markedly higher than before the 1960s, with 1973, 1980, 1988, 1991, 1998 and 2005 outstanding in this respect. There has been a marked decrease in the frequency of frost across Australia. According to the Bureau of Meteorology, Australia's annual mean temperature for 2009 was 0.9 °C above the 1961–90 average, making it the nation's second-warmest year since high-quality records began in 1910. According to the Bureau of Meteorology's 2011 Australian Climate Statement, Australia had lower than average temperatures in 2011 as a consequence of a La Niña weather pattern; however, "the country's 10-year average continues to demonstrate the rising trend in temperatures, with 2002–2011 likely to rank in the top two warmest 10-year periods on record for Australia, at above the long-term average". Furthermore, 2014 was Australia's third warmest year since national temperature observations commenced in 1910. Sea level rise The Australian Government released a report saying that up to 247,600 houses are at risk from flooding from a sea level rise of 1.1 metres. There were 39,000 buildings located within 110 metres of 'soft' erodible shorelines, at risk from a faster erosion due to sea level rise. Adaptive responses to this specific climate change threat are often incorporated in the coastal planning policies and recommendations at the state level. For instance, the Western Australia State Coastal Planning Policy established a sea level rise benchmark for initiatives that address the problem over a 100-year period. Lower projections indicate that sea levels will rise by 40 to 90 cm upon the end of the century Water (droughts and floods) Bureau of Meteorology records since the 1860s show that a 'severe' drought has occurred in Australia, on average, once every 18 years. Australia is already the driest populated continent in the world. Rainfall in southwestern Australia has decreased by 10–20% since the 1970s, while southeastern Australia has also experienced a moderate decline since the 1990s. Rainfall is expected to become heavier and more infrequent, as well as more common in summer rather than in winter. In June 2008 it became known that an expert panel had warned of long-term, maybe irreversible, severe ecological damage for the whole Murray-Darling basin if it did not receive sufficient water by October of that year. Water restrictions were in place in many regions and cities of Australia in response to chronic shortages resulting from the 2008 drought. In 2004 paleontologist Tim Flannery predicted that unless it made drastic changes the city of Perth, Western Australia, could become the world's first ghost metropolis—an abandoned city with no more water to sustain its population. In 2019 the Drought and Water Resources Minister of Australia David Littleproud, said, that he "totally accepts" the link between climate change and drought in Australia because he "lives it". He called for a reduction in greenhouse gas emission and massive installation of renewable energy. Former leader of the Nationals Barnaby Joyce said that if the drought became more fierce and dams were not built, the Coalition risks "political annihilation". According to the 2022 IPCC report, there has been an increase in flooding episodes and other catastrophic weather events because of global warming. These unusual weather changes in include rainfall in the north and severe droughts in the south. Less rainfall means less streamflow of water for major cities. The IPCC recommends a step up to our adaptation and finance policies in our systems to keep up with the drastic impacts of climate change for a sustainable development. Water resources Healthy and diverse vegetation is essential to river health and quality, and many of Australia's most important catchments are covered by native forest, maintaining a healthy ecosystem. Climate change will affect growth, species composition and pest incursion of native species and in turn, will profoundly affect water supply from these catchments. Increased re-afforestation in cleared catchments also has the prospect for water losses. Between 1970 and 2024, 28% of Australia's Hydrological Reference Stations showed a significant decrease in streamflow while 4% showed a significant increase. The stations with increases were all in northern Australia while those with decreases were largely in southern Australia. The CSIRO predicted that the additional effects in Australia of a temperature rise of between only 1 and 2 °C will be: 12–25% reduction inflow in the Murray River and Darling River basin. 7–35% reduction in Melbourne's water supply. Bushfires There is an increase in fire activity in Australia since 1950. The causes include "more dangerous fire weather conditions, increased risk factors associated with pyroconvection, including fire-generated thunderstorms, and increased ignitions from dry lightning, all associated to varying degrees with anthropogenic climate change". Firefighting officials are concerned that the effects of climate change will increase the frequency and intensity of bushfires under even a "low global warming" scenario. A 2006 report, prepared by CSIRO Marine and Atmospheric Research, Bushfire CRC, and the Australian Bureau of Meteorology, identified South Eastern Australia as one of the three most fire-prone areas in the world, and concluded that an increase in fire-weather risk is likely at most sites over the next several decades, including the average number of days when the McArthur Forest Fire Danger Index rating is very high or extreme. It also found that the combined frequencies of days with very high and extreme FFDI ratings are likely to increase 4–25% by 2020 and 15–70% by 2050, and that the increase in fire-weather risk is generally largest inland. Former Australian Greens leader, Bob Brown said that the fires were "a sobering reminder of the need for this nation and the whole world to act and put at a priority the need to tackle climate change". The Black Saturday Royal Commission recommended that "the amount of fuel-reduction burning done on public land each year should be more than doubled". In 2018, the fire season in Australia began in the winter. August 2018 was hotter and windier than the average. Those meteorological conditions led to a drought in New South Wales. The Government of the state already gave more than $1 billion to help the farmers. The hotter and drier climate led to more fires. The fire seasons in Australia are lengthening and fire events became more frequent in the latest 30 years. These trends are probably linked to climate change. The 2019–20 Australian bushfire season was by some measures Australia's "worst bushfire season on record". In New South Wales, the fires burnt through more land than any other blazes in the past 25 years, in addition to being the state's worst bushfire season on record. NSW also experienced the longest continuously burning bushfire complex in Australia's history, having burnt more than , with flames being reported. Approximately 3 billion animals were killed or displaced by the bushfires and this made them one of the worst natural disasters in recorded history. The chance of reaching the climatic conditions that fuels the fires became more than four times bigger since the year 1900 and will become eight times more likely to occur if the temperature will rise by 2 degrees from the preindustrial level. In December 2019 the New South Wales Government declared a state of emergency after record-breaking temperatures and prolonged drought exacerbated the bushfires. In 2019 bushfires linked to climate change created air pollution 11 times higher that the hazardous level in many areas of New South Wales. Many medical groups called to protect people from "public health emergency" and moving on from fossil fuels. According to the United Nations Environment Programme the megafires in Australia in 2019–2020 are probably linked to climate change that created the unusually dry and hot weather conditions. This is part of a global trend. Brazil, the United States, the Russian Federation, and the Democratic Republic of the Congo all face similar problems. By the second week of January the fires burned a territory of approximately 100,000 square kilometres close to the territory of England, killed one billion animals and caused large economic damage. Researchers claim that the exceptionally strong wildfires in 2019–2020 were impossible without the effects of climate change. More than one-fifth of Australian forests were burned in one season, which was completely unprecedented. They say that: "In the case of recent events in Australia, there is no doubt that the record temperatures of the past year would not be possible without anthropogenic influence, and that under a scenario where emissions continue to grow, such a year would be average by 2040 and exceptionally cool by 2060." Climate change probably also caused drier weather conditions in Australia by impacting Indian Ocean Dipole, which also increase fires. In average, below 2% of Australian forests burn annually. Climate change has increased the likelihood of the wildfires in 2019–2020 by at least 30%, but researchers said the result is probably conservative. Extreme weather events Rainfall patterns and the degree of droughts and storms brought about by extreme weather conditions are likely to be affected. The CSIRO predicts that a temperature rise of between 2 and 3 °C on the Australian continent could incur some of the following extreme weather occurrences, in addition to standard patterns: Wind speeds of tropical cyclones could intensify by 5 to 10%. In 100 years, strong tides would increase by 12–16% along eastern Victoria's coast. The forest fire danger indices in New South Wales and Western Australia would grow by 10% and the forest fire danger indices in south, central and north-east Australia would increase by more than 10%. Heatwaves A report in 2014 revealed that, due to the change in climatic patterns, heat waves were found to be increasingly more frequent and severe, with an earlier start to the season and longer duration. Since temperatures began to be recorded in 1910, they have increased by an average of 1 °C, with most of this change occurring from 1950 onwards. This period has seen the frequency and intensity of extreme heat events increase. Summer 2013–14 was warmer than average for the entirety of Australia. Both Victoria and South Australia saw record-breaking temperatures. Adelaide recorded a total of 13 days reaching 40 °C or more, 11 of which reached 42 °C or more, as well as its fifth-hottest day on record—45.1 °C on 14 January. The number of days over 40 °C beat the previous record of summer 1897–1898, when 11 days above 40 °C were recorded. Melbourne recorded six days over 40 °C, while nighttime temperatures were much warmer than usual, with some nights failing to drop below 30 °C. Overall, the summer of 2013–2014 was the third-hottest on record for Victoria, fifth-warmest on record for New South Wales, and sixth-warmest on record for South Australia. This heatwave has been directly linked to climate change, which is unusual for specific weather events. Following the 2014 event, it was predicted that temperatures might increase by up to 1.5 °C by 2030. 2015 was Australia's fifth-hottest year on record, continuing the trend of record-breaking high temperatures across the country. According to Australian Climate Council in 2017 Australia had its warmest winter on record, in terms of average maximum temperatures, reaching nearly 2 °C above average. January 2019 was the hottest month ever in Australia with average temperatures exceeding . Ecosystems and biodiversity Sustained climate change could have drastic effects on the ecosystems of Australia. For example, rising ocean temperatures and continual erosion of the coasts from higher water levels will cause further bleaching of the Great Barrier Reef. Beyond that, Australia's climate will become even harsher, with more powerful tropical cyclones and longer droughts. The Department of Climate Change said in its Climate Change Impacts and Costs fact sheet: "...ecologically rich sites, such as the Great Barrier Reef, Queensland Wet Tropics, Kakadu Wetlands, Australian Alpine areas, south-western Australia and sub- Antarctic islands are all at risk, with significant loss of biodiversity projected to occur by 2020". It also said: "Very conservatively, 90 Australian animal species have so far been identified at risk from climate change, including mammals, insects, birds, reptiles, fish, and amphibians from all parts of Australia." Australia has some of the world's most diverse ecosystems and natural habitats, and it may be this variety that makes them the Earth's most fragile and at-risk when exposed to climate change. The Great Barrier Reef is a prime example. Over the past 20 years it has experienced unparalleled rates of bleaching. Additional warming of 1 °C is expected to cause substantial losses of species and of associated coral communities. The CSIRO predicts that the additional results in Australia of a temperature rise of between 2 and 3 °C will be: 97% of the Great Barrier Reef bleached annually. 10–40% loss of principal habitat for Victoria and montane tropical vertebrate species. 92% decrease in butterfly species' primary habitats. 98% reduction in Bowerbird habitat in Northern Australia. 80% loss of freshwater wetlands in Kakadu (30 cm sea level rise). A study conducted in 2024, suggests that worsening climate scenarios may have significant impacts on the habitat area for vertebrates and vascular plant species in Australia. This data suggests that in 2030, the habitat area among species appears consistent across climate scenarios, however, by 2090, a significant change is predicted, as deteriorating climate conditions are associated with reductions in habitat area for biodiversity. Great Barrier Reef The Great Barrier Reef could be killed as a result of the rise in water temperature forecast by the IPCC. A UNESCO World Heritage Site, the reef has experienced unprecedented rates of bleaching over the past two decades, and additional warming of only 1 °C is anticipated to cause considerable losses or contractions of species associated with coral communities. Lord Howe Island The coral reefs of the World Heritage-listed Lord Howe Island could be killed as a result of the rise in water temperature forecast by the IPCC. As of April 2019, approximately 5% of the coral is dead. Impacts on people Economic impacts According to the Climate Commission (now the Climate Council) report in 2013, the extreme heatwaves, flooding and bushfires striking Australia have been intensified by climate change and will get worse in future in terms of their impacts on people, property, communities and the environment. The summer of 2012/2013 included the hottest summer, hottest month and hottest day on record. The cost of the 2009 bushfires in Victoria was estimated at A$4.4bn (£3bn) and the Queensland floods of 2010/2011 cost over A$5bn. In 2008 the Treasurer and the Minister for Climate Change and Water released a report that concluded the economy will grow with an emissions trading scheme in place. A report released in October 2009 by the Standing Committee on Climate Change, Water, Environment and the Arts, studying the effects of a 1-metre sea level rise, quite possible within the next 30–60 years, concluded that around 700,000 properties around Australia, including 80,000 buildings, would be inundated, the collective value of these properties is estimated at $155 billion. In 2019 the Australian Bureau of Agricultural and Resource Economics and Sciences published a report about the impact of climate change on the profitability of the Australian agriculture, saying that the profit of the Australian farms was cut by 22% due to climate change in the years 2000–2019. According to the 2022 IPCC report Australia will lose billions of dollars due to loss of life, and physical damages. These natural disasters are caused by climate change and increasing global warming will worsen these events. The report estimates that under 2 degrees of warming Australia will lose $115 billion in the next decade, and $350 billion in the next twenty years. If warming goes up to under 3 degrees of warming Australia's economy will lose $200 billion and $600 billion by 2042. Agriculture forestry and livestock Small changes caused by global warming, such as a longer growing season, a more temperate climate and increased concentrations, may benefit Australian crop agriculture and forestry in the short term. However, such benefits are unlikely to be sustained with increasingly severe effects of global warming. Changes in precipitation and consequent water management problems will further exacerbate Australia's current water availability and quality challenges, both for commercial and residential use. The CSIRO predicts that the additional results in Australia of a temperature rise of between 3 and 4 °C will be: 32% possibility of diminished wheat production (without adaptation). 45% probability of wheat crop value being beneath present levels (without adaptation). 55% of primary habitat lost for Eucalyptus. 25–50% rise in common timber yield in cool and wet parts of South Australia. 25–50% reduction in common timber yield in North Queensland and the Top End. 6% decrease in Australian net primary production (for 20% precipitation decrease) 128% increase in tick-associated losses in net cattle production weight. Electricity demand Use of domestic air conditioners during severe heatwaves can double electricity demand, placing great stress on electricity generation and transmission networks, and lead to load shedding. In addition, bushfires can damage electricity lines, while repairing power poles and power line damages is often restricted during hot and dry weather because of high fire risks. Impacts on housing Settlements and infrastructure Global warming could lead to substantial alterations in climate extremes, such as tropical cyclones, heat waves and severe precipitation events. This would degrade infrastructure and raise costs through intensified energy demands, maintenance for damaged transportation infrastructure, and disasters, such as coastal flooding. In the coastal zone, sea level rise and storm surge may be more critical drivers of these changes than either temperature or precipitation. The CSIRO describes the additional impact on settlements and infrastructure for rises in temperature of only 1 to 2 °C: A 22% rise in 100-year storm surge height around Cairns; as a result, the area flooded doubles. Human settlements Climate change will have a higher impact on Australia's coastal communities, due to the concentration of population, commerce and industry. Climate modelling suggests that a temperature rise of 1–2 °C will result in more intense storm winds, including those from tropical cyclones. Combine this with sea level rise, and the result is greater flooding, due to higher levels of storm surge and wind speed. The impact of climate change on insurance against catastrophes. Proceedings of Living with Climate Change Conference. Canberra, 19 December.) Tourism of coastal areas may also be affected by coastal inundation and beach erosion, as a result of sea level rise and storm events. At higher levels of warming, coastal impacts become more severe with higher storm winds and sea levels. Property A report released in October 2009 by the Standing Committee on Climate Change, Water, Environment and the arts, studying the effects of a 1-metre sea level rise, possible within the next 30–60 years, concluded that around 700,000 properties around Australia, including 80,000 buildings, would be inundated. The collective value of these properties is estimated at $150 billion. A 1-metre sea level rise would have massive impacts, not just on property and associated economic systems, but in displacement of human populations throughout the continent. Queensland is the state most at risk due to the presence of valuable beachfront housing. Impacts on foreign policy and national security Several prominent reports and decision makers are concerned that climate change affects Australia’s national security. A 2023 assessment of the Australian Office of National Intelligence on the security implications of climate change (commissioned by Anthony Albanese) remains classified. These concerns are tied to broader debates about climate security. Climate change is unlikely to trigger large-scale migration movements to Australia. Research shows that people adversely affected by climate change often lack the resources to migrate over large distances (they are adapt in-situ or move to nearby places). This is particularly the case for Australia, which is an island nation with tough border controls and immigration policies. Climate change can cause major challenges to Australia’s foreign policies. Pacific Island countries, which are highly vulnerable to climate change, have repeatedly blamed Australia for not being active enough in mitigating climate change. With geopolitical tensions between Australia and China on the rise, these countries are of high relevance for the Australian government. In addition, existing data suggest that several countries in Australia’s neighbourhood (e.g., Indonesia, the Philippines, large parts of South Asia) and some key partner governments (e.g., India, Papua New Guinea) are very vulnerable to climate-related unrest and conflict. Several reports also warn that climate change poses significant challenges to the capacities of the Australian Defence Force. Many military bases are located close to the coastline, which is threatened by sea-level rise and more intense storms. Civilian infrastructure relevant to military operations (like transports networks and power lines) is also adversely affected by climate change, for instance when floods wash away key supply roads. More extreme heat days complicate military training and put a heavier toll on equipment, particularly in northern Australia. Finally, the Australian Defence Force will be called upon more often to provide disaster relief within the country and internationally, further straining its resources. Health impacts The CSIRO predicts that the additional results in Australia of a temperature rise of between only 1 and 2 °C will be: Southward spread of malaria receptive zones. Risk of dengue fever among Australians increases from 170,000 people to 0.75–1.6 million. 10% increase in diarrhoeal diseases among Aboriginal children in central Australia. 100% increase in a number of people exposed to flooding in Australia. Increased influx of refugees from the Pacific Islands. Based on some predictions for 2070, data suggests that people who are not accustomed to the warmer climate may experience as much as 45 days per year where they are unable to tolerate being outside, compared to the current 4–6 days per year. Impacts on indigenous Australians Indigenous Australians have a millennia long history of responding and adapting to social and environmental changes. Indigenous Australians have a high level of situated traditional knowledge and historical knowledge about climate change. However, the exposure of Indigenous Australians to climate change impacts is exacerbated by existing socio-economic disadvantages which are linked to colonial and post-colonial marginalisation. Some of these changes include a rise in sea levels, getting hotter and for a longer period of time, and more severe cyclones during the cyclone season. Climate issues include wild fires, heatwaves, floods, cyclones, rising sea levels, rising temperatures, and erosion. The communities most affected by climate changes are those in the North where Aboriginal and Torres Strait Islander people make up 30% of the population. Aboriginal Australians and Torres Strait Islander communities located in the coastal north are the most disadvantaged due to social and economic issues and their reliance on traditional land for food, culture, and health. This has begged the question for many community members in these regions, should they move away from this area or remain present. Many Aboriginal people live in rural and remote agricultural areas across Australia, especially in the Northern and Southern areas of the continent. There are a variety of different climate impacts on different Aboriginal communities which includes cyclones in the Northern region and flooding in Central Australia which negatively impacts cultural sites and therefore the relationship between indigenous people and the places that hold their traditional knowledge. Other effects include sea level rise, loss of land and hunting ground, changes in fire regimes, increased severity and duration of wet and dry seasons as well as reduced numbers of animals in the sea, rivers and creeks. Vulnerability The vulnerability comes from remote location where indigenous groups live, lower socio-economic status, and reliance of natural systems for economic needs. Disadvantages which are compounding Indigenous peoples vulnerability to climate change include inadequate health and educational services, limited employment opportunities as well as insufficient infrastructure. Top down institutions have also restricted Indigenous Australians ability to contribute to climate policy frameworks and have their culture and practices recognised. Many of the economic, political, and social-ecological issues present in indigenous communities are long term effects from colonialism and the continued marginalization of these communities. These issues are aggravated by climate change and environmental changes in their respective regions. Indigenous people are seen as particularly vulnerable to climate change because they already live in poverty, poor housing and have poor educational and health services, other socio-political factors place them at risk for climate change impacts. Indigenous people have been portrayed as victims and as vulnerable populations for many years by the media. Aboriginal Australians believe that they have always been able to adapt to climate changes in their geographic areas. Many communities have argued for more community input into strategies and ways to adapt to climate issues instead of top down approaches to combating issues surrounding environmental change. This includes self-determination and agency when deciding how to respond to climate change including proactive actions. Indigenous people have also commented on the need to maintain their physical and mental well-being in order to adapt to climate change which can be helped through the kinship relationships between community members and the land they occupy. In Australia, Aboriginal people have argued that in order for the government to combat climate change, their voices must be included in policy making and governance over traditional land. Much of the government and institutional policies related to climate change and environmental issues in Australia has been done so through a top down approach. Indigenous communities have stated that this limits and ignores Aboriginal Australian voices and approaches. Due to traditional knowledge held by these communities and elders within those communities, traditional ecological knowledge and frameworks are necessary to combat these and a variety of different environmental issues. Heat and drought Fires and droughts in Australia, including the Northern regions, occur mostly in savannas because of current environmental changes in the region. The majority of the fire prone areas in the savanna region are owned by Aboriginal Australian communities, the traditional stewards of the land. Aboriginal Australians have traditional landscape management methods including burning and clearing the savanna areas which are the most susceptible to fires. Traditional landscape management declined in the 19th century as Western landscape management took over. Today, traditional landscape management has been revitalized by Aboriginal Australians, including elders. This traditional landscape practices include the use of clearing and burning to get rid of old growth. Though the way in which indigenous communities in this region manage the landscape has been banned, Aboriginal Australian communities who use these traditional methods actually help in reducing greenhouse gas emissions. Impact of climate change on health Increased temperatures, wildfires, and drought are major issues in regard to the health of Aboriginal Australian communities. Heat poses a major risk to elderly members of communities in the North. This includes issues such as heat stroke and heat exhaustion. Many of the rural indigenous communities have faced thermal stress and increased issues surrounding access to water resources and ecological landscapes. This impacts the relationship between Aboriginal Australians and biodiversity, as well as impacts social and cultural aspects of society. Aboriginal Australians who live in isolated and remote traditional territories are more sensitive than non-indigenous Australians to changes that effect the ecosystems they are a part of. This is in large part due to the connection that exists between their health (including physical and mental), the health of their land, and the continued practice of traditional cultural customs. Aboriginal Australians have a unique and important relationship with the traditional land of their ancestors. Because of this connection, the dangerous consequences of climate change in Australia has resulted in a decline in health including mental health among an already vulnerable population. In order to combat health disparities among these populations, community based projects and culturally relevant mental and physical health programs are necessary and should include community members when running these programs. Traditional knowledge Indigenous people have always responded and adapted to climate change, including indigenous people of Australia. Aboriginal Australian people have existed in Australia for tens of thousands of years. Due to this continual habitation, Aboriginal Australians have observed and adapted to climatic and environmental changes for millennia which uniquely positions them to be able to respond to current climate changes. Though these communities have shifted and changed their practices overtime, traditional ecological knowledge exists that can benefit local and indigenous communities today. This knowledge is part of traditional cultural and spiritual practices within these indigenous communities. The practices are directly tied to the unique relationship between Aboriginal Australians and their ecological landscapes. This relationship results in a socio-ecological system of balance between humans and nature Indigenous communities in Australia have specific generational traditional knowledge about weather patterns, environmental changes and climatic changes. These communities have adapted to climate change in the past and have knowledge that non-Indigenous people may be able to utilize to adapt to climate change currently and in the future. Indigenous people have not been offered many opportunities or provided with sufficient platforms to influence and contribute their traditional knowledge to the creation of current international and local policies associated to climate change adaptation. Although, Indigenous people have pushed back on this reality, by creating their own platforms and trying to be active members in the conversation surrounding climate change including at international meetings. Specifically, Indigenous people of Australia have traditional knowledge to adapt to increased pressures of global environmental change. Though some of this traditional knowledge was not utilised and conceivably lost with the introduction of white settlers in the 18th century, recently communities have begun to revitalize these traditional practices. Australian Aboriginal traditional knowledge includes language, cultural, spiritual practices, mythology and land management. Responses to climate change Indigenous knowledge has been passed down through the generations with the practice of oral tradition. Given the historical relationship between the land and the people and the larger ecosystem Aboriginal Australians choose to stay and adapt in similar ways to their ancestors before them. Aboriginal Australians have observed short and long term environmental changes and are highly aware of weather and climate changes. Recently, elders have begun to be utilised by indigenous and non-indigenous communities to understand traditional knowledge related to land management. This includes seasonal knowledge means indigenous knowledge pertaining to weather, seasonal cycles of plants and animals, and land and landscape management. The seasonal knowledge allows indigenous communities to combat environmental changes and may result in healthier social-ecological systems. Much of traditional landscape and land management includes keeping the diversity of flora and fauna as traditional foodways. Ecological calendars is one traditional framework used by Aboriginal Australian communities. These ecological calendars are way for indigenous communities to organize and communicate traditional ecological knowledge. The ecological calendars includes seasonal weather cycles related to biological, cultural, and spiritual ways of life. Mitigation Climate change mitigation focuses on steps taken to reduce greenhouse gas emissions. It is the set of preventative measures taken to curb global warming and climate change. Examples would be investing in clean fuel and using renewable energy such as wind and solar power. According to the CSIRO and Garnaut Climate Change Review, climate change is expected to have numerous adverse effects on many species, regions, activities and much infrastructure and areas of the economy and public health in Australia. The Stern Report and Garnaut Review on balance expect these to outweigh the costs of mitigation. The World Resources Institute identifies policy uncertainty and over-reliance on international markets as the top threats to Australia's GHG mitigation. Emissions reductions Internationally, Australia pledged as part of Paris Agreement to reduce emissions by 43% by 2030 and achieve net zero emissions by 2050. Domestically, the Clean Energy Act 2011 addresses GHG with an emissions cap, carbon price, and subsidies. Emissions by the electric sector are addressed by Renewable Energy targets at multiple scales, Australian Renewable Energy Agency (ARENA), Clean Energy Finance Corporation (CEFC), carbon capture and storage flagships, and feed-in tariffs on solar panels. Emissions by the industrial sector are addressed by the Energy Efficiency Opportunities (EEO) program. Emissions by the building sector are addressed by building codes, minimum energy performance standards, Commercial Building Disclosure program, state energy-saving obligations, and the National Energy Saving Initiative. Emissions by the transportation sector are addressed by reduced fuel tax credits and vehicle emissions performance standards. Emissions by the agricultural sector are addressed by the Carbon Farming Initiative and state land-clearing laws. Emissions by the land use sector are addressed by the Clean Energy Future Package, which consists of the Carbon Farming Futures program, Diversity Fund, Regional Natural Resources Management Planning for Climate Change Fund, Indigenous Carbon Farming Fund, and Carbon Farming Skills program. Forestry and forest-related options for carbon sinks In Australia, forestry and forest-related options are the most significant and most easily achieved carbon sink making up 105 Mt per year CO2-e or about 75 per cent of the total figure attainable for the Australian state of Queensland from 2010 to 2050. Among the forestry options, forestry with the primary aim of carbon storage (called carbon forestry) has the highest attainable carbon storage capacity (77 Mt CO2-e/yr) while strategy balanced with biodiversity plantings can return 7–12 times more native vegetation for a 10%–30% reduction of carbon storage performance. Legal strategies to encourage this form of biosequestration include permanent protection of forests in National Parks or on the World Heritage List, properly funded management and bans on use of rainforest timbers and inefficient uses such as woodchipping old growth forest. Policies and legislation to achieve mitigation Paris Agreement The Paris agreement is a legally international agreement adopted at the COP 21, its main goal is to limit global warming to below 1.5 °C, compared to pre-industrial levels. The Nationally Determined Contributions (NDCs) are the plans to fight climate change adapted for each country. Every party in the agreement has different goals based on its own historical climate records and country's circumstances. All the goals for each country are stated in their NDC. Australia's target regarding reductions from 2005 year levels: 26–28% reduction of greenhouse gases (GHG) until 2030 from 2005 levels. In 2022 the new Australian government officially declared the update of the targets to 43% reduction by 2030 and net zero emissions by 2050. Gases covered in reductions: Carbon dioxide (), Methane (), Nitrous oxide (), Hydrofluorocarbon (HFCs), Perfluorinated compound (PFCs), Sulfur hexafluoride (SF6) and Nitrogen trifluoride (NF3). Countries have different ways to achieve the established goals depending on resources. Australia's developed approach to support the NDC climate change plan is the following: Enabling new technologies with low emissions and promoting economic growth. Establish regional hydrogen exports to strengthen the country's industry and fund research in the field and enable distribution. Improve charging and refueling infrastructure to enable companies and fleets to integrate new more sustainable vehicle technology. The country has created a development fund whose purpose is for projects concerning carbon dioxide capture. The fund is for storage, use and carbon capture. Investments in technological development that reduces emissions in the sectors of agriculture, industry, transport and manufacturing. Climate solution package to increase investment in projects to generate clean energy. The package also includes extra funds to support development in the hard-to-reach sectors. Australia has a legalised obligation for the major emitting sectors in the country where the emissions are to be kept below their baseline. Australia has through funds such as Australia emission reduction fund contributed with 60 million tonnes reduction of greenhouse gases. The fund enables businesses to earn carbon credits. This is done by storing or preventing emissions through new sustainable techniques. State legislation Victoria The Climate Change Act was adopted in 2017 and is part of a broader Victorian environmental legislation taking climate change into account. It establishes a net-zero emission target by 2050 and interim targets set every five years to adapt and keep Victoria on track with the 2050 goal. Adaptation According to the IPCC's 2001 Assessment Report, no matter how much effort is put into mitigating climate change, some amount of climate change cannot be avoided. The report shared that climate change adaptation should complement mitigation efforts. Adaptation is the approach that focuses on alleviating current problems brought about by global warming and climate change. It is the attempt to live with the changes in the environment and the economy that global warming has generated and will continue to generate. In short, it involves taking action to deal with the problems brought about by global warming and climate change. Examples include building better flood defences and avoiding the building of residential areas near low-lying, flood-prone areas. In cities with a proven vulnerability to climate change, investment is likely to require the strengthening of urban infrastructure, including storm drain systems, water supply and treatment plants, and protection or relocation of solid waste management and power generation facilities. Coastal regions are likely to need large investment in physical infrastructure projects, specifically projects related to the effects of rising sea levels. Projects such as the construction of protective barriers against rising sea levels, the building of dams to retain and manage water, the redesign and development of port facilities and the improvement of the defence systems at coastal areas should be carried out. Federal, state and territory policy makers have supported a National Biodiversity and Climate Change Action Plan that works to adapt to the impacts of climatic change and manage the effects on wildlife. National government programs Regional natural resource management (NRM) organisations Federal natural resource goals, government agencies and non-government organizations established 56 regional natural resource management (NRM) organisations beginning in the mid-1990s. NRM organisations fall under the federal government Natural Heritage Trust. NRM operate according to individual constitutions, usually by the state government and others by community associations. Their boards are appointed by either the local government or community stakeholders. NRM Planning for Climate Fund, put $13.6 million toward helping NRMs plan land use in light of climate change by building a base of detailed climatic information. National Climate Change Adaptation Programme The Minister for Energy and Emissions Reduction has come up with the National Climate Change Adaptation Programme which aims to work with industries, scientific organisations, residents and other governments to create workable solutions. Some A$14 million over a period of four years (2008–2012) is to be spent on this initiative. The programme has forged strong research links in at-risk areas such as the Great Barrier Reef. Research conducted in the Great Barrier Reef is focused on developing methods to deal with climate change to protect the reef. It is hoped that this work will create a universal model for sustainable, cost-effective reef development. According to the programme's brochure: "National greenhouse mitigation policies and programmes are projected to reduce emissions by 94 million tonnes by 2010 – the equivalent of removing every motor vehicle in Australia from the road! However, the greenhouse gases already in the atmosphere and the growing emissions from around the world will affect our climate. Adaptation to climate change will complement action to reduce greenhouse gases". Climate Adaptation Flagship The Australian Commonwealth Scientific and Industrial Research Organisation (CSIRO) started the Climate Adaptation Flagship. Its aim is "enabling Australia to adapt more effectively to the impacts of climate change and variability and informing national planning, regulation and investment decisions". This is part of the National Research Flagships Program. It is designed to bring various stakeholders, i.e. research companies, industries, international connections, eminent scientists and CSIRO, together in hope of delivering practical solutions that address the pressing issues of Australia. The Climate Adaptation Flagship project concerns both climate variability (or non-human causes, as defined by the United Nations Framework Convention on Climate Change) and climate change. The research budget for this Flagship for the year 2008–09 is close to A$30 million. There are four research prongs to this project: Pathways to adaptation; Sustainable cities and costs; Managing species and natural ecosystems; Adaptive primary industries, enterprises and communities. National Climate Change Adaptation Research Facility The National Climate Change Adaptation Research Facility (NCCARF) is hosted by Griffith University in Queensland and "leads the research community in a national interdisciplinary effort to generate the information needed by decision-makers in government and in vulnerable sectors and communities to manage the risks of climate change impacts". The key roles of NCCARF include: developing National Adaptation Research Plans to identify critical gaps in the information available to decision-makers synthesising existing and emerging national and international research on climate change impacts and adaptation and developing targeted communication products undertaking a program of Integrative Research to address national priorities, and establishing and maintaining Adaptation Research Networks to link together key researchers and assist them in focusing on national research priorities. The facility is a partnership between the Australian government's Department of Climate Change and Griffith University, with a consortium of funding partners and universities drawn from across the country. The Local Adaptations Pathway Program The Australian government is of the view that local government is critical in managing the impacts of climate change and seeks to assist local councils in studying and applying adaptation options. The programme is the Australian government's initiative to enable councils to go through climate change risk assessments and come up with action plans to prepare for the impacts the phenomenon may have on local society. Up to A$50,000 will be released. A list of councils successful in procuring the funding is provided on the programme's website. Policies and legislation In November 1981, the Office of National Assessments (intelligence agency) presented prime minister Malcolm Fraser with a classified-confidential assessment noting scientific acceptance of the greenhouse effect and resultant "measurably warmer" temperatures and "related climatic changes", and also projecting effects of possible doubling and quadrupling of atmospheric levels by the middle and end of the 21st century. The assessment focused on the implications for the country's fossil fuel industry. In the late 1980s and early 1990s, there was clear Australian consensus about the need for action on climate change between the two major political parties. However, following the 1991 recession, incoming right wing governments began framing science of climate change as a continuing debate. In 1997, Australia joined the United States as the only countries to not ratify the Kyoto Protocol. With voters influenced by events like the Millennium drought and 2006 film An Inconvenient Truth, both parties went to the 2007 election promising action on climate change, with the then opposition calling climate change the "greatest moral, economic and social challenge of our time". The incumbent Howard government lost, and the incoming Labor government immediately ratified the Kyoto Protocol. In 2009, before a bill could be passed, with the support of opposition leader Malcolm Turnbull, the opposition changed leaders to Tony Abbott, and supported by The Greens but for the opposite reason that Rudd's scheme was too weak and potentially locked in failure, blocked Rudd's Carbon Pollution Reduction Scheme. In 2010, the Rudd government decided to delay the implementation of the Carbon Pollution Reduction Scheme (CPRS) until the end of the first commitment period of the Kyoto Protocol (ending in 2012). They cited the lack of bipartisan support for the CPRS and slow international progress on climate action as the reasons for the decision. In turn, the delay was strongly criticised by the Federal Opposition as well as community and grassroots action groups such as GetUp. Following the unsuccessful Copenhagen Summit, the Rudd was replaced by Gillard as prime minister, who stated that "there will be no 'carbon tax' under the government I lead". The Gillard Labor government established several government entities to manage Australia's response to climate change: The Climate Change Authority, an independent statutory body that provides advice and performs research for the federal government on climate change. The Clean Energy Regulator, an independent statutory body that administers federal government schemes to measure and reduce Australia's greenhouse gas emissions. The Australian Renewable Energy Agency (ARENA), a corporate body that manages renewable energy programs. The Clean Energy Finance Corporation (CEFC), a government-owned corporation that invests in clean energy technologies. In 2011, Parliament passed the Clean Energy Act 2011, which introduced carbon pricing in Australia, colloquially known as a 'carbon tax'. It required large businesses, defined as those emitting over 25,000 tons of carbon dioxide equivalent annually, to purchase emissions permits. The strong backlash led by opposition leader Abbott led to her being replaced as leader by Rudd, then Abbott at the next election. Under his leadership, Australia became the first country to repeal a carbon pricing program. In 2015, Abbott was replaced as prime minister by Minister for Communications Malcolm Turnbull under the condition that his climate policy would not change. Australia attended the 2015 United Nations Climate Change Conference and adopted the Paris Agreement. In limiting further action on climate change, Australia joined Russia, Turkey and Brazil in citing US President Trump's promise to withdraw from the Paris Agreement. In 2018, Turnbull was replaced by Scott Morrison as leader of the Liberal Party and prime minister. Morrison won the 2019 election with an unchanged climate policy. In June 2021, the Sustainable Development Report 2021 scored Australia last out of 193 United Nations member countries for action taken to reduce global greenhouse gas emissions, scoring 10 out of 100 in an assessment of fossil fuel emissions, emissions associated with imports and exports, and policies for pricing carbon. In May 2022, the Coalition lost the federal election to the Labor Party, led by Anthony Albanese. In a machinery of government change, a new Department of Climate Change, Energy, the Environment and Water will be established. The new government committed to a 43% reduction in Australia's emissions by 2030 (compared to 2005 levels), and net zero emissions by 2050. History of climate change policy in Australia Domestic action to address climate change in Australia began in 1989, when Senator Graham Richardson proposed the first greenhouse gas emission reduction target of 20% by 2005. The Australian Government rejected this target. In 1990, Ros Kelly and Jon Kerin announced that the Australian Government would adhere to the goals initially proposed by Richardson but not to any economic detriment. Australia signed the UNFCCC in 1992. This was followed by the release of the National Greenhouse Response Strategy (NGRS), which provided states and territories with the mechanisms to adhere to UNFCCC emission guidelines. Australia attended the first session of the Conference of the Parties to the UNFCCC in Berlin in March 1995. Throughout the 1990s, Australia regularly failed to meet its own emission targets and those set by the UNFCCC. In 1997, Prime Minister John Howard announced that by 2010, an additional 2% of electricity would be sustainably sourced. The following year, the Australian Greenhouse Office (AGO) was established to monitor greenhouse gas reductions. The AGO later combined with the Department of Environment and Heritage. In April 1998, Australia became a party to the Kyoto Declaration. The Declaration was ratified in 2007 under Prime Minister Kevin Rudd. In the Renewable Energy (Electricity) Act 2000, the Federal Government introduced the Mandatory Renewable Energy Target program, which aimed to sustainably source 10% of electrical energy by 2010. In 2011, the Mandatory Renewable Energy Target program was divided into the Large-Scale Renewable Energy Target and the Small-Scale Renewable Energy Scheme. In January 2003, the New South Wales State Government implemented the Greenhouse Gas Reduction Scheme (GGRS), which allowed carbon emissions to be traded. Under Rudd, the Labor Government proposed the Carbon Pollution Reduction Scheme, which was intended to take effect in 2010. This scheme was rejected by the Greens for being too permissive and by Tony Abbott's Coalition for being economically detrimental. Under Prime Minister Julia Gillard, the Labor Party passed the Clean Energy Act 2011 to establish a carbon tax and put a price on greenhouse gas emissions. This carbon tax was a divisive partisan issue. In 2012, the Coalition ran a campaign to repeal the carbon tax. Upon election victory in September 2013, Prime Minister Tony Abbott passed the Clean Energy Legislation (Carbon Tax Repeal) Bill. In replacement of the carbon tax, Abbott introduced the Direct Action Scheme to financially reward businesses for voluntarily reducing their carbon emissions. This was followed by a decision not to participate in the 19th session of the Conference of the Parties to the UNFCCC (COP-19). Australia became a party to the Paris Agreement in 2015. In the agreement, Australia committed to reducing its emissions by 26% by 2030. In 2019, Prime Minister Scott Morrison was criticised for a lack of commitment to addressing climate change while taking a vacation during the 2019 bushfires. International cooperation Internationally, Australia contributed to the creation of the Asia Pacific Rain Forest Partnership, International Coral Reef Initiative, International Partnership for Blue Carbon, Mission Innovation, Clean Energy Ministerial Forum, International Solar Alliance, and the Kigali Amendment to the Montreal Protocol. The government has also provided $1 billion to assist developing countries in reducing GHG emissions, partly through the United Nations Framework Convention on Climate Change Green Climate Fund. Australia's scientists also provide data on climate, emissions, impacts, and mitigation options for the Intergovernmental Panel on Climate Change assessments. Under the Paris Agreement, Australia has committed to reducing emission by 26-28% below 2005 levels. This would mean reducing emissions by half per capita and by two-thirds across the economy. The Department of Environment and Energy noted in a 2017 review that no one policy could achieve what multiple, sector-specific ones have. This approach has manifested in Australia meeting its first Kyoto Protocol target. Australia is now bound to reducing emissions to at least 5% by 2020 under the Copenhagen Accord and Cancun Agreements and 0.5% less than 1990 levels by 2020 under their second target for the Kyoto Protocol. While Australia opposed a 1.5 °C target at the 2015 negotiations for the Paris Agreement, in 2019, they supported the Kainaki II Declaration of the Pacific Islands Forum, which included this target. In 2022, Australia discussed hosting COP29 with its Pacific island neighbours in 2024 at the Pacific Islands Forum. In November 2023 it was announced that Australia will offer 280 Tuvalu citizens displaced by climate change permanent residency in Australia per year, as part of a broad bilateral deal. Society and culture Politics Despite the support of a clear scientific consensus, climate change has been a divisive or controversial issue in Australian politics since the 2000s. It has sometimes been referred to as a "culture war" in the country. Conservatives have generally opposed climate mitigation policies and renewable energy, instead favouring or supporting the country's coal and fossil fuels industries, which make up a large part of the economy. Proposed carbon pricing during the premiership of Julia Gillard proved highly divisive, and was later repealed under Tony Abbott. Climate change was a key issue in the 2022 federal election, where the Australian Labor Party and teal independents made gains in part due to promoting environmental policies. Australian conservatives, with the support of strongly climate-skeptical media, have long opposed climate change mitigation and changes to energy policy. This is partly a strategy to foster the support of the country's coal and the fossil fuel industry, which are highly influential and a large employer in the country. Activism Climate change protests have taken place in Australia during the 21st century.In 2005, with support from Uniting Church and Catholic Earthcare, the Australian Conservation Foundation and the National Council of Churches Australia produced a brochure, Changing Climate, Changing Creation, which was distributed to churches across the country to call for action on climate change. Rising Tide held environmental direct action protests in February 2007, where more than 100 small and medium-sized craft, including swimmers and people on surfboards, gathered in Newcastle harbour. Young people from the Real Action On Climate Change shut down two coal-fired power stations in September 2007. A 2009 "Walk Against Warming" drew 40,000 participants in Melbourne. The Say Yes demonstrations took place on 5 June 2011, in which 45,000 people demonstrated in every major city nationwide in support of carbon pricing policies. Thousands of Australian children took part in school strikes for climate in 2018 and 2019. The September 2019 climate strikes attracted an estimated 180,000 to 300,000 participants across eight Australian capital cities and 140 urban centres, making it one of the largest protests in the country's history and one of the largest climate protests globally. Approximately 2,500 businesses also took part. The response to the 2019–20 Australian bushfire season sparked protests in Sydney, Canberra, Melbourne, Victoria, Brisbane, Hobart, and outside the Australian High Commission in London. Prime Minister Scott Morrison was criticised for climate denial in the wake of the bushfires. Extinction Rebellion held rallies in London, Berlin, Madrid, Copenhagen and Stockholm calling for stronger climate action. Direct action group Blockade Australia began disruptive activism in 2021 and 2022. In 2023 another Rising Tide water blockade was held in Newcastle during which 3000 people took part and 109 were arrested. Despite the introduction of tougher penalties in New South Wales for such activity the majority of those facing court received dismissals with no conviction with magistrates acknowledging the protesters as “valuable contributors to society" and commending their “muscular good character. Litigation Groups including Rising Tide and Queensland Conservation have initiated legal challenges to coal mines under the Commonwealth EPBC legislation. In late 2006, Queensland Conservation lodged an objection to the greenhouse gas emissions from a large coal mine expansion proposed by Xstrata Coal Queensland Pty Ltd. QC's action aimed to have the true costs of the greenhouse gas emissions from coal mining recognised. The Newlands Coal Mine Expansion will produce 28.5 million tonnes of coal over its fifteen years of operation. The mining, transport and use of this coal will emit 84 million tonnes of into the atmosphere. Queensland Conservation aims to have reasonable and practical measures imposed on new mines to avoid, reduce or offset the emissions from the mining, transport and use of their coal. The Land and Resources Tribunal ruled against the case. Media coverage Projected impacts by location The impacts of climate change will vary significantly across Australia. The Australian Government appointed Climate Commission have prepared summary reports on the likely impacts of climate change for regions across Australia, including: Queensland, NSW, Victoria and Tasmania. Capital cities Adelaide Adelaide will get hotter and drier with rainfall predicted to decline 8% to 29% by 2090 and average temperature to increase between 4 and 0.9 degrees. The number of days above 35 degrees will increase by 50% in 2090 and the number of days above 40 degrees will double. Bringing it close to Northampton, Western Australia, for temperature and Kadina, South Australia, for rainfall. Sea levels will rise with predictions between 39 and 61 cm by 2090. And extreme seas are predicted to rise as well, with the CSIRO predicting buildings in Port Adelaide would need to be raised by 50 to 81 cm to keep the amount of flooding incidents the same as recorded between 1986 and 2005. Brisbane In a RCP 4.5 scenario Brisbane's temperature will be similar to that of Rockhampton today while rainfall will be closest to Gympie. The CSIRO predicts rainfall in Brisbane will fall between -23% (235 mm) and -4% (45.3 mm) annually by 2090 while temperature will rise between 4.2° and 0.9°. The number of hot days and hot nights will double by 2050, with many people needing to avoid outdoor activity in summer. Further urban growth increases the number of hot nights even further. Hot nights increase deaths amongst the elderly. Rainfall will be deposited in less frequent more intense rain events, fire days will also get more frequent while frost days will decrease. Sea levels are predicted to rise by 80 cm by 2100 and there will be more frequent sea level extremes. Darwin In a RCP 4.5 scenario Darwin's temperature will be similar to that of Daly River now, with its rainfall most like that of Milikapiti. In a RCP 8.5 scenario, indicating higher greenhouse gas emissions, Darwin's temperature loses any close comparison in Australia being significantly hotter than every town in Australia is today (with the exclusion of Halls Creek in Autumn). Sydney Suburbs of Sydney like Manly, Botany, Narrabeen, Port Botany, and Rockdale, which lie on rivers like the Parramatta, face risks of flooding in low-lying areas such as parks (like Timbrell Park and Majors Bay Reserve), or massive expenses in rebuilding seawalls to higher levels. Sea levels are predicted to rise between 38 and 66 cm by 2090. Temperature in Sydney will increase between 0.9° and 4.2°, while rainfall will decrease between -23% and -4% by 2090. Bringing Sydney's climate close to that of Beaudesert today (under a RCP 8.5 scenario). Different parts of Sydney will warm differently with the greatest impact expected in Western Sydney and Hawkesbury, these areas can expect 5 to 10 additional hot days by 2030. Similarly future rainfall patterns will be different to those today, with more rain expected to fall in summer and autumn and less expected in Winter and Spring. Fire danger days will increase in number by 2070. Melbourne Sea levels are projected to rise between 0.37 cm and 0.59 cm at Williamstown (the closest covered point) by 2090. At the higher end of this scale areas in and around Melbourne would be impacted. With some of the most vulnerable areas being the Docklands development and several marinas and berths in Port Phillip. Melbourne's climate will become similar in terms of total rainfall and average temperature to that of Dubbo today, with temperatures warming between 0.9° and 3.8° and total annual rainfall falling between -10% and -4% by 2090. Rainfall patterns will also change with 20% less rainfall predicted during spring in 2050, which may impact the severity of summer bushfires. The increases in temperature and decrease in rainfall will have a series of follow on effects on the city, including a possible 35% reduction in trees in Melbourne by 2040. And more frequent ambulance callouts and more deaths due to heatwaves. Climate change will cost Melbourne City $12.6bn by 2050 and be closer to Wangaratta's climate. Perth In 2090 Perth is predicted to have the rainfall of Yanchep today and the temperature of Geraldton using the RCP 4.5 scenario. Rainfall is predicted to fall between -29% (-226 mm) and -8% (-66 mm) and temperature predicted to rise between 0.9° and 4°. Perth may see the number of days above 35° increase from 28 per year on average to 36 in 2030, and to between 40 and 63 in 2090. While frost days will decrease. Rainfall will increase in intensity while decreasing on average. Drought days in the south west as a whole may increase by as much as 80% versus 20% for Australia. The danger from fire will increase with more fire days for all of Western Australia. Hobart By 2090 Hobart's climate will warm between 3.8° and 0.9°, rainfall will decline between 4% and 10%. The temperature pattern will be similar to Port Lincoln while rainfall will be closer to Condoblin's today in a RCP 8.5 scenario. Warm spells are likely to last longer and rainfall will trend to more intense rain events dumping less rain annually, increasing the risk of erosion and flooding. Flooding on the Derwent river will become more regular and extreme with a current 1-in-100-year event being possibly a 2-to-6-year event in 2090. Hobart's fire season will get longer. States Victoria By 2050, Victoria's annual temperature will increase up to 2.4 °C, with twice the number of very hot days compared to 1986-2005, longer fire seasons, less rainfall and snowfall in cool season and a rise in sea levels about 24 cm. Historical aspects Pre-instrumental climate change Paleoclimatic records indicate that during glacial maxima Australia was extremely arid, with plant pollen fossils showing deserts as far as northern Tasmania and a vast area of less than 12% vegetation cover over all of South Australia and adjacent regions of other states. Forest cover was largely limited to sheltered areas of the east coast and the extreme southwest of Western Australia. During these glacial maxima the climate was also much colder and windier than today. Minimum temperatures in winter in the centre of the continent were as much as lower than they are today. Hydrological evidence for dryness during glacial maxima can also be seen at major lakes in Victoria's Western District, which dried up between around 20,000 and 15,000 years ago and re-filled from around 12,000 years ago. During the early Holocene, there is evidence from Lake Frome in South Australia and Lake Woods near Tennant Creek that the climate between 8,000 and 9,500 years ago and again from 7,000 to 4,200 years ago was considerably wetter than over the period of instrumental recording since about 1885. The research that gave these records also suggested that the rainfall flooding Frome was certainly summer-dominant rainfall because of pollen counts from grass species. Other sources suggest that the Southern Oscillation may have been weaker during the early Holocene and rainfall over northern Australia less variable as well as higher. The onset of modern conditions with periodic wet season failure is dated at around 4,000 years before the present. In southern Victoria, there is evidence for generally wet conditions except for a much drier spell between about 3,000 and 2,100 years before the present, when it is believed Lake Corangamite fell to levels well below those observed between European settlement and the 1990s. After this dry period, Western District lakes returned to their previous levels fairly quickly and by 1800 they were at their highest levels in the forty thousand years of record available. Elsewhere, data for most of the Holocene are deficient, largely because methods used elsewhere to determine past climates (like tree-ring data) cannot be used in Australia owing to the character of its soils and climate. Recently, however, coral cores have been used to examine rainfall over those areas of Queensland draining into the Great Barrier Reef. The results do not provide conclusive evidence of man-made climate change, but do suggest the following: There has been a marked increase in the frequency of very wet years in Queensland since the end of the Little Ice Age, a theory supported by there being no evidence for any large Lake Eyre filling during the LIA. The dry era of the 1920s and 1930s may well have been the driest period in Australia over the past four centuries. A similar study, not yet published, is planned for coral reefs in Western Australia. Records exist of floods in a number of rivers, such as the Hawkesbury, from the time of first settlement. These suggest that, for the period beginning with the first European settlement, the first thirty-five years or so were wet and were followed by a much drier period up to the mid-1860s, when usable instrumental records started. Development of an instrumental network for climate records Although rain gauges were installed privately by some of the earliest settlers, the first instrumental climate records in Australia were not compiled until 1840 at Port Macquarie. Rain gauges were gradually installed at other major centres across the continent, with the present gauges in Melbourne and Sydney dating from 1858 and 1859, respectively. In eastern Australia, where the continent's first large-scale agriculture began, a large number of rain gauges were installed during the 1860s and by 1875 a comprehensive network had been developed in the "settled" areas of that state. With the spread of the pastoral industry to the north of the continent during this period, rain gauges were established extensively in newly settled areas, reaching Darwin by 1869, Alice Springs by 1874, and the Kimberley, Channel Country and Gulf Savannah by 1880. By 1885, most of Australia had a network of rainfall reporting stations adequate to give a good picture of climatic variability over the continent. The exceptions were remote areas of western Tasmania, the extreme southwest of Western Australia, Cape York Peninsula, the northern Kimberley and the deserts of northwestern South Australia and southeastern Western Australia. In these areas good-quality climatic data were not available for quite some time after that. Temperature measurements, although made at major population centres from days of the earliest rain gauges, were generally not established when rain gauges spread to more remote locations during the 1870s and 1880s. Although they gradually caught up in number with rain gauges, many places which have had rainfall data for over 125 years have only a few decades of temperature records.	Physical sciences	Climate change	Earth science
11664252	https://en.wikipedia.org/wiki/Dromiidae	Dromiidae	Dromiidae is a family of crabs, often referred to as sponge crabs. They are small or medium-sized crabs which get their name from the ability to shape a living sponge into a portable shelter for themselves. A sponge crab cuts out a fragment from a sponge and trims it to its own shape using its claws. The last two pairs of legs are shorter than other legs and bend upward over the crab's carapace, to hold the sponge in place. The sponge grows along with the crab, providing a consistent shelter. Subfamilies and genera The family Dromiidae contains the following subfamilies and genera: Dromiinae Alainodromia Ameridromia Ascidiophilus Austrodromidia Barnardomia Conchoecetes Costadromia Cryptodromia Cryptodromiopsis Desmodromia Dromia Dromidia Dromidiopsis Dromilites Epigodromia Epipedodromia Eudromidia Exodromidia Foredromia Fultodromia Haledromia Hemisphaerodromia Homalodromia Kerepesia Kromtitis Lamarckdromia Lauridromia Lewindromia Lucanthonisia Mclaydromia Metadromia Moreiradromia Noetlingia Paradromia Petalomera Platydromia Pseudodromia Speodromia Stebbingdromia Sternodromia Stimdromia Takedromia Tumidodromia Tunedromia Hypoconchinae Hypoconcha Sphaerodromiinae Eodromia Frodromia Sphaerodromia	Biology and health sciences	Crabs and hermit crabs	Animals
11669530	https://en.wikipedia.org/wiki/Nose	Nose	A nose is a sensory organ and respiratory structure in vertebrates. It consists of a nasal cavity inside the head, and an external nose on the face. The external nose houses the nostrils, or nares, a pair of tubes providing airflow through the nose for respiration. Where the nostrils pass through the nasal cavity they widen, are known as nasal fossae, and contain turbinates and olfactory mucosa. The nasal cavity also connects to the paranasal sinuses (dead-end air cavities for pressure buffering and humidification). From the nasal cavity, the nostrils continue into the pharynx, a switch track valve connecting the respiratory and digestive systems. In humans, the nose is located centrally on the face and serves as an alternative respiratory passage especially during suckling for infants. The protruding nose that is completely separate from the mouth part is a characteristic found only in therian mammals. It has been theorized that this unique mammalian nose evolved from the anterior part of the upper jaw of the reptilian-like ancestors (synapsids). Air treatment Acting as the first interface between the external environment and an animal's delicate internal lungs, a nose conditions incoming air, both as a function of thermal regulation and filtration during respiration, as well as enabling the sensory perception of smell. Hair inside nostrils filter incoming air, as a first line of defense against dust particles, smoke, and other potential obstructions that would otherwise inhibit respiration, and as a kind of filter against airborne illness. In addition to acting as a filter, mucus produced within the nose supplements the body's effort to maintain temperature, as well as contributes moisture to integral components of the respiratory system. Capillary structures of the nose warm and humidify air entering the body; later, this role in retaining moisture enables conditions for alveoli to properly exchange O2 for CO2 (i.e., respiration) within the lungs. During exhalation, the capillaries then aid recovery of some moisture, mostly as a function of thermal regulation, again. Sense of direction The wet nose of dogs is useful for the perception of direction. The sensitive cold receptors in the skin detect the place where the nose is cooled the most and this is the direction a particular smell that the animal just picked up comes from. Structure in air-breathing forms In amphibians and lungfish, the nostrils open into small sacs that, in turn, open into the forward roof of the mouth through the choanae. These sacs contain a small amount of olfactory epithelium, which, in the case of caecilians, also lines a number of neighbouring tentacles. Despite the general similarity in structure to those of amphibians, the nostrils of lungfish are not used in respiration, since these animals breathe through their mouths. Amphibians also have a vomeronasal organ, lined by olfactory epithelium, but, unlike those of amniotes, this is generally a simple sac that, except in salamanders, has little connection with the rest of the nasal system. In reptiles, the nasal chamber is generally larger, with the choanae located much further back in the roof of the mouth. In crocodilians, the chamber is exceptionally long, helping the animal to breathe while partially submerged. The reptilian nasal chamber is divided into three parts: an anterior vestibule, the main olfactory chamber, and a posterior nasopharynx. The olfactory chamber is lined by olfactory epithelium on its upper surface and possesses a number of turbinates to increase the sensory area. The vomeronasal organ is well-developed in lizards and snakes, in which it no longer connects with the nasal cavity, opening directly into the roof of the mouth. It is smaller in turtles, in which it retains its original nasal connection, and is absent in adult crocodilians. Birds have a similar nose to reptiles, with the nostrils located at the upper rear part of the beak. Since they generally have a poor sense of smell, the olfactory chamber is small, although it does contain three turbinates, which sometimes have a complex structure similar to that of mammals. In many birds, including doves and fowls, the nostrils are covered by a horny protective shield. The vomeronasal organ of birds is either under-developed or altogether absent, depending on the species. The nasal cavities in mammals are both fused into one. Among most species, they are exceptionally large, typically occupying up to half the length of the skull. In some groups, however, including primates, bats, and cetaceans, the nose has been secondarily reduced, and these animals consequently have a relatively poor sense of smell. The nasal cavity of mammals has been enlarged, in part, by the development of a palate cutting off the entire upper surface of the original oral cavity, which consequently becomes part of the nose, leaving the palate as the new roof of the mouth. The enlarged nasal cavity contains complex turbinates forming coiled scroll-like shapes that help to warm the air before it reaches the lungs. The cavity also extends into neighbouring skull bones, forming additional air cavities known as paranasal sinuses. In cetaceans, the nose has been reduced to one or two blowholes, which are the nostrils that have migrated to the top of the head. This adaptation gave cetaceans a more streamlined body shape and the ability to breathe while mostly submerged. Conversely, the elephant's nose has elaborated into a long, muscular, manipulative organ called the trunk. The vomeronasal organ of mammals is generally similar to that of reptiles. In most species, it is located in the floor of the nasal cavity, and opens into the mouth via two nasopalatine ducts running through the palate, but it opens directly into the nose in many rodents. It is, however, lost in bats, and in many primates, including humans. In fish Fish have a relatively good sense of smell. Unlike that of tetrapods, the nose has no connection with the mouth, nor any role in respiration. Instead, it generally consists of a pair of small pouches located behind the nostrils at the front or sides of the head. In many cases, each of the nostrils is divided into two by a fold of skin, allowing water to flow into the nose through one side and out through the other. The pouches are lined by olfactory epithelium, and commonly include a series of internal folds to increase the surface area, often forming an elaborate "olfactory rosette". In some teleosts, the pouches branch off into additional sinus-like cavities, while in coelacanths, they form a series of tubes. In the earliest vertebrates, there was only one nostril and olfactory pouch, and the nasal passage was connected to the hypophysis. The same anatomy is observed in the most primitive living vertebrates, the lampreys and hagfish. In gnathostome ancestors, the olfactory apparatus gradually became paired (presumably to allow sense of direction of smells), and freeing the midline from the nasal passage allowed evolution of jaws.	Biology and health sciences	Nervous system	null
4046824	https://en.wikipedia.org/wiki/Gauss%E2%80%93Seidel%20method	Gauss–Seidel method	In numerical linear algebra, the Gauss–Seidel method, also known as the Liebmann method or the method of successive displacement, is an iterative method used to solve a system of linear equations. It is named after the German mathematicians Carl Friedrich Gauss and Philipp Ludwig von Seidel. Though it can be applied to any matrix with non-zero elements on the diagonals, convergence is only guaranteed if the matrix is either strictly diagonally dominant, or symmetric and positive definite. It was only mentioned in a private letter from Gauss to his student Gerling in 1823. A publication was not delivered before 1874 by Seidel. Description Let be a square system of linear equations, where: When and are known, and is unknown, the Gauss–Seidel method can be used to iteratively approximate . The vector denotes the initial guess for , often for . Denote by the -th approximation or iteration of , and by the approximation of at the next (or -th) iteration. Matrix-based formula The solution is obtained iteratively via where the matrix is decomposed into a lower triangular component , and a strictly upper triangular component such that . More specifically, the decomposition of into and is given by: Why the matrix-based formula works The system of linear equations may be rewritten as: The Gauss–Seidel method now solves the left hand side of this expression for , using the previous value for on the right hand side. Analytically, this may be written as Element-based formula However, by taking advantage of the triangular form of , the elements of can be computed sequentially for each row using forward substitution: Notice that the formula uses two summations per iteration which can be expressed as one summation that uses the most recently calculated iteration of . The procedure is generally continued until the changes made by an iteration are below some tolerance, such as a sufficiently small residual. Discussion The element-wise formula for the Gauss–Seidel method is related to that of the (iterative) Jacobi method, with an important difference: In Gauss-Seidel, the computation of uses the elements of that have already been computed, and only the elements of that have not been computed in the -th iteration. This means that, unlike the Jacobi method, only one storage vector is required as elements can be overwritten as they are computed, which can be advantageous for very large problems. However, unlike the Jacobi method, the computations for each element are generally much harder to implement in parallel, since they can have a very long critical path, and are thus most feasible for sparse matrices. Furthermore, the values at each iteration are dependent on the order of the original equations. Gauss-Seidel is the same as successive over-relaxation with . Convergence The convergence properties of the Gauss–Seidel method are dependent on the matrix . Namely, the procedure is known to converge if either: is symmetric positive-definite, or is strictly or irreducibly diagonally dominant. The Gauss–Seidel method may converge even if these conditions are not satisfied. Golub and Van Loan give a theorem for an algorithm that splits into two parts. Suppose is nonsingular. Let be the spectral radius of . Then the iterates defined by converge to for any starting vector if is nonsingular and . Algorithm Since elements can be overwritten as they are computed in this algorithm, only one storage vector is needed, and vector indexing is omitted. The algorithm goes as follows: algorithm Gauss–Seidel method is inputs: , repeat until convergence for from 1 until do for from 1 until do if ≠ then end if end (-loop) end (-loop) check if convergence is reached end (repeat) Examples An example for the matrix version A linear system shown as is given by: Use the equation in the form where: Decompose into the sum of a lower triangular component and a strict upper triangular component : The inverse of is: Now find: With and the vectors can be obtained iteratively. First of all, choose , for example The closer the guess to the final solution, the fewer iterations the algorithm will need. Then calculate: As expected, the algorithm converges to the solution: . In fact, the matrix is strictly diagonally dominant, but not positive definite. Another example for the matrix version Another linear system shown as is given by: Use the equation in the form where: Decompose into the sum of a lower triangular component and a strict upper triangular component : The inverse of is: Now find: With and the vectors can be obtained iteratively. First of all, we have to choose , for example Then calculate: In a test for convergence we find that the algorithm diverges. In fact, the matrix is neither diagonally dominant nor positive definite. Then, convergence to the exact solution is not guaranteed and, in this case, will not occur. An example for the equation version Suppose given equations and a starting point . At any step in a Gauss-Seidel iteration, solve the first equation for in terms of ; then solve the second equation for in terms of just found and the remaining ; and continue to . Then, repeat iterations until convergence is achieved, or break if the divergence in the solutions start to diverge beyond a predefined level. Consider an example: Solving for and gives: Suppose is the initial approximation, then the first approximate solution is given by: Using the approximations obtained, the iterative procedure is repeated until the desired accuracy has been reached. The following are the approximated solutions after four iterations. The exact solution of the system is . An example using Python and NumPy The following iterative procedure produces the solution vector of a linear system of equations: import numpy as np ITERATION_LIMIT = 1000 # initialize the matrix A = np.array( [ [10.0, -1.0, 2.0, 0.0], [-1.0, 11.0, -1.0, 3.0], [2.0, -1.0, 10.0, -1.0], [0.0, 3.0, -1.0, 8.0], ] ) # initialize the RHS vector b = np.array([6.0, 25.0, -11.0, 15.0]) print("System of equations:") for i in range(A.shape[0]): row = [f"{A[i,j]:3g}x{j+1}" for j in range(A.shape[1])] print("[{0}] = [{1:3g}]".format(" + ".join(row), b[i])) x = np.zeros_like(b, np.float_) for it_count in range(1, ITERATION_LIMIT): x_new = np.zeros_like(x, dtype=np.float_) print(f"Iteration {it_count}: {x}") for i in range(A.shape[0]): s1 = np.dot(A[i, :i], x_new[:i]) s2 = np.dot(A[i, i + 1 :], x[i + 1 :]) x_new[i] = (b[i] - s1 - s2) / A[i, i] if np.allclose(x, x_new, rtol=1e-8): break x = x_new print(f"Solution: {x}") error = np.dot(A, x) - b print(f"Error: {error}") Produces the output: System of equations: [ 10x1 + -1x2 + 2x3 + 0x4] = [ 6] [ -1x1 + 11x2 + -1x3 + 3x4] = [ 25] [ 2x1 + -1x2 + 10x3 + -1x4] = [-11] [ 0x1 + 3x2 + -1x3 + 8x4] = [ 15] Iteration 1: [ 0. 0. 0. 0.] Iteration 2: [ 0.6 2.32727273 -0.98727273 0.87886364] Iteration 3: [ 1.03018182 2.03693802 -1.0144562 0.98434122] Iteration 4: [ 1.00658504 2.00355502 -1.00252738 0.99835095] Iteration 5: [ 1.00086098 2.00029825 -1.00030728 0.99984975] Iteration 6: [ 1.00009128 2.00002134 -1.00003115 0.9999881 ] Iteration 7: [ 1.00000836 2.00000117 -1.00000275 0.99999922] Iteration 8: [ 1.00000067 2.00000002 -1.00000021 0.99999996] Iteration 9: [ 1.00000004 1.99999999 -1.00000001 1. ] Iteration 10: [ 1. 2. -1. 1.] Solution: [ 1. 2. -1. 1.] Error: [ 2.06480930e-08 -1.25551054e-08 3.61417563e-11 0.00000000e+00] Program to solve arbitrary number of equations using Matlab The following code uses the formula function x = gauss_seidel(A, b, x, iters) for i = 1:iters for j = 1:size(A,1) x(j) = (b(j) - sum(A(j,:)'.x) + A(j,j)*x(j)) / A(j,j); end end end	Mathematics	Linear algebra	null
4047104	https://en.wikipedia.org/wiki/Jacobi%20method	Jacobi method	In numerical linear algebra, the Jacobi method (a.k.a. the Jacobi iteration method) is an iterative algorithm for determining the solutions of a strictly diagonally dominant system of linear equations. Each diagonal element is solved for, and an approximate value is plugged in. The process is then iterated until it converges. This algorithm is a stripped-down version of the Jacobi transformation method of matrix diagonalization. The method is named after Carl Gustav Jacob Jacobi. Description Let be a square system of n linear equations, where: When and are known, and is unknown, we can use the Jacobi method to approximate . The vector denotes our initial guess for (often for ). We denote as the k-th approximation or iteration of , and is the next (or k+1) iteration of . Matrix-based formula Then A can be decomposed into a diagonal component D, a lower triangular part L and an upper triangular part U:The solution is then obtained iteratively via Element-based formula The element-based formula for each row is thus:The computation of requires each element in except itself. Unlike the Gauss–Seidel method, we cannot overwrite with , as that value will be needed by the rest of the computation. The minimum amount of storage is two vectors of size n. Algorithm Input: , (diagonal dominant) matrix A, right-hand side vector b, convergence criterion Output: Comments: pseudocode based on the element-based formula above while convergence not reached do for i := 1 step until n do for j := 1 step until n do if j ≠ i then end end end increment k end Convergence The standard convergence condition (for any iterative method) is when the spectral radius of the iteration matrix is less than 1: A sufficient (but not necessary) condition for the method to converge is that the matrix A is strictly or irreducibly diagonally dominant. Strict row diagonal dominance means that for each row, the absolute value of the diagonal term is greater than the sum of absolute values of other terms: The Jacobi method sometimes converges even if these conditions are not satisfied. Note that the Jacobi method does not converge for every symmetric positive-definite matrix. For example, Examples Example question A linear system of the form with initial estimate is given by We use the equation , described above, to estimate . First, we rewrite the equation in a more convenient form , where and . From the known values we determine as Further, is found as With and calculated, we estimate as : The next iteration yields This process is repeated until convergence (i.e., until is small). The solution after 25 iterations is Example question 2 Suppose we are given the following linear system: If we choose as the initial approximation, then the first approximate solution is given by Using the approximations obtained, the iterative procedure is repeated until the desired accuracy has been reached. The following are the approximated solutions after five iterations. The exact solution of the system is . Python example import numpy as np ITERATION_LIMIT = 1000 # initialize the matrix A = np.array([[10., -1., 2., 0.], [-1., 11., -1., 3.], [2., -1., 10., -1.], [0.0, 3., -1., 8.]]) # initialize the RHS vector b = np.array([6., 25., -11., 15.]) # prints the system print("System:") for i in range(A.shape[0]): row = [f"{A[i, j]}*x{j + 1}" for j in range(A.shape[1])] print(f'{" + ".join(row)} = {b[i]}') print() x = np.zeros_like(b) for it_count in range(ITERATION_LIMIT): if it_count != 0: print(f"Iteration {it_count}: {x}") x_new = np.zeros_like(x) for i in range(A.shape[0]): s1 = np.dot(A[i, :i], x[:i]) s2 = np.dot(A[i, i + 1:], x[i + 1:]) x_new[i] = (b[i] - s1 - s2) / A[i, i] if x_new[i] == x_new[i-1]: break if np.allclose(x, x_new, atol=1e-10, rtol=0.): break x = x_new print("Solution: ") print(x) error = np.dot(A, x) - b print("Error:") print(error) Weighted Jacobi method The weighted Jacobi iteration uses a parameter to compute the iteration as with being the usual choice. From the relation , this may also be expressed as . Convergence in the symmetric positive definite case In case that the system matrix is of symmetric positive-definite type one can show convergence. Let be the iteration matrix. Then, convergence is guaranteed for where is the maximal eigenvalue. The spectral radius can be minimized for a particular choice of as follows where is the matrix condition number.	Mathematics	Linear algebra	null
4047117	https://en.wikipedia.org/wiki/Motorcycle%20taxi	Motorcycle taxi	A motorcycle taxi, or cart bike or bike taxi, is a licensed form of transport in some countries. The taxi typically carries one passenger, who "rides pillion" behind the motorcycle operator. Multiple passengers are common in some countries. Brazil According to some sources, motorcycle taxi service in Brazil began in 1994, in Crateús, in the state of Ceará, when an employee of the Bank of Brazil Other sources state that it started in Bauru, São Paulo, in 1996, when an unemployed biker hung a banner across the road to the city, reading "help a biker racing to 1.00 real." Today, almost all Brazilian cities have motorcycle taxi services. Recently, they have appeared in poorer and less urban areas, where young people increasingly support themselves by driving them. Typically, the fare is a flat fee, regardless of the distance traveled. However, the charge may vary according to the time of day or day of the week, or increase for distances that are greater than usual. Licensing requirements for motorcycle taxis vary by municipality. Small towns tend not to regulate them at all, while in larger cities, they are regulated in much the same way as taxicabs. In July 2009, the Brazilian Senate approved standards for motorcycle taxi drivers and motorcycle couriers. They must be at least 21 years old, have held a Category A drivers licence for at least two years, and have attended a training course. Cambodia In Phnom Penh and other cities in Cambodia, motorcycle taxis are widely available as a form of low-cost public transport. Motorcycle taxi drivers, who are almost exclusively male, are called motodops (). They tend to hang around outside major tourist attractions, office buildings, public markets, and near the corners of residential streets. There is no regulated system of training or bike maintenance and no common uniform, so anyone on his way home from the market might offer you a ride (and the driver's intentions can generally be trusted, the state of his bike, a little less so). Always negotiate the fare in advance (use gestures, if necessary). Don't expect a motodop to understand English or to read a map - he'll likely flag somebody down who can help translate or navigate, if necessary. Fares vary depending on distance and weather but should always be cheaper than a tuk-tuk. Fares are higher at night and when embarking from tourist areas. You'll get a better rate if you can negotiate in passable Khmer, but have a heart: these are generally the folks that live on a few dollars or less per day. As of 2014, helmet laws apply only to drivers, so bring your own helmet if you're worried about safety, but it's not legally required. The omnipresent ‘moto’ is the most common and fastest form of public transportation. Motos can be found virtually everywhere in town, just step to the curb and they will find you. Motos cost from 1500R-4000R for a trip in town and $6-$8 per day. Prices go up at night and for multiple passengers. Cameroon Motorcycle taxis are also the most common form of transportation in Maroua, Cameroon. Multiple passengers are carried on most trips; as many as four children are sometimes carried on a single motorcycle. Helmets are rarely used, but the traffic and speed are moderate in the city. Short distances cost about 200 francs, less than US$1. China In mainland China, motorcycle taxi service can be traced back to the late 1980s and early 1990s. There are currently motorcycle taxis throughout China, including in Beijing, Shanghai, and Guangzhou. They are popular primarily due to their low cost: the fare for short trips is just 5 yuan (less than US$1) per person. India In Goa, India, motorcycle taxis are required to be licensed. Driven by men called pilots, they are much cheaper than other taxis, although a passenger can only carry a backpack as luggage. In some parts of the state, motorcycle drivers are legally required to wear helmets, but any passengers riding pillion are not. Motorcycle taxis can usually be identified by their distinctive yellow and black colours. There is a practice to fix the fare in advance, and trips are not metered. In last few years, a few companies such as Rapido, Uber and Ola have come up in multiple cities in India providing bike taxi services. With the Central Government's rule of allowing two-wheelers as legal and commercial vehicles and 8 states already legalized the same, it has become easier for the companies to design a working framework to provide easy and comfortable commute to the people. Indonesia Motorcycle taxis are a very common form of unlicensed transport in Indonesia, where they are known as ojek. Ojek can be found in most areas of the country, from towns where traffic jams commonly hinder other forms of transport, to rural areas inaccessible by four-wheeled vehicles. Because of the traffic, ojek are often the fastest form of transport, especially in big urban areas such as Jakarta, Surabaya, or Medan. Many people choose them over taxicabs, which are safer, but slower and more expensive. Many ojek drivers either own their vehicles or are buying them on credit, although in some areas, stolen motorcycles are common. The widespread availability of cheap, domestic motorcycles made by Honda, Yamaha, and Suzuki, and even cheaper ones imported from China, as well as credit schemes with which to purchase these, have resulted in the rapid growth of ojek. The ease with which driver's licences can be obtained has also been a contributing factor. Before the trip begins, the passenger usually haggles with the driver over the fare, which generally ranges from 5,000 to 10,000 rupiah (about US$0.50 to US$1.00) for short trips, longer trips will be more expensive. The fare may be different from one city to another city, as big city ojek will have higher fares than the smaller city ojek. Indonesia traffic law requires motorcycle riders and passengers to wear helmets; often on ojek, however, only the driver does so. Although the driver will sometimes provide a helmet for the passenger, more often, drivers simply avoid main streets, and the attention of police. The name of Gojek is derived from the word ojek. Mexico In Mexico, there are thousands of motorcycle taxis. Their arrangements are informal (not traditional companies). They have precarious working conditions, long hours (11.3 hours a day), low wages (US$59.18 per week), and no social protections or benefits. 6.3% reported suffering from a disease, 49.5% corresponded to musculoskeletal conditions and only 11.6% were affiliated to any health system. 53.8% are owners of the vehicle and, although it does not seem to influence physical illness (P=0.03), it is related to the psychosocial ones (P=0.260). Nigeria Nigeria has about three million motorcycle taxis, locally called okadas, with over one million in Lagos alone. In Lagos, new rules prohibit okadas from carrying pregnant women or children. Authorities say okadas will be stopped from driving the wrong way, and the number of roads on which they are authorized to travel will be sharply reduced. Philippines Motorcycle taxis in the Philippines usually have sidecars, or seats extended sideways, often by use of a T-shaped crossbeam. The latter type of taxi is known as habal-habal or a skylab, owing to its crude resemblance to the Skylab space station which orbited the Earth in the 1970s. Covered, three-wheel auto rickshaws, known as tricycles, are also a common mode of transport. Angkas is a Philippine motorcycle vehicle for hire and delivery service company based in Makati, Metro Manila. Its competition in passenger market is JoyRide. Motorcycle taxis were deemed illegal in 2020 due to possible exposure of passengers and riders to COVID-19 when in contact with each other, especially in the cities. Thailand Motorcycle taxis (, ; , ; or , ) are a common form of public transport in Bangkok and most other cities, towns, and villages in Thailand. They are generally used for short trips. In Bangkok, there are motorcycle taxi queues on many sois, or side streets, and the queues are regulated by land transport authority. Licensed motorcycle taxi operators wear orange vests with yellow number plates. The driving license with photo and driver's details in form of yellow card is placed on the back of the driver where the passenger can see clearly. In compliance with Thailand's motorcycle helmet law, many (but not all) drivers carry a spare helmet to offer to passengers. Bangkok locals generally only use motorcycle taxis when they need to get somewhere fast, as metered taxi-cabs can not only be more expensive for short trips but also slower than flat-rate motorcycles. Therefore, motorcycle taxi-drivers in Bangkok have built their reputation on delivering service as quickly as possible and tend to drive very fast and weave through traffic. United Kingdom Motorcycle taxi service in London began in 1990 as a niche industry. All equipment is provided for the passenger, along with an intercom system linking the rider and passenger. The motorcycles have racks that can hold a carry-on suitcase, for trips to local airports, especially Stansted, Gatwick, and City. The bikes are now licensed by Transport for London and the Public Carriage Office, which also license London's black cabs. United States Moto Limos Club, a motorcycle for-hire service, started in California and New York City in 2011. As of 2012, the business filing was not renewed and of 2015 the filing was considered suspended. Passengers were not able to hail the motorcycles on the street; instead, a yearly individual or corporate membership fee is charged, plus an hourly rate. Experienced riders, many former Police motorcycle riders, carried clients on Honda Gold Wings, and in California, can bypass traffic congestion by lane splitting. Passengers were provided with helmets, airbag vests, and in-helmet, Bluetooth cell phones. The service also bought several Can-Am Spyders, before realizing they were not capable of splitting lanes. Vietnam Nimble motorcycle taxis, which surpass buses in speed and mobility, comprise one of the most popular modes of transportation in Vietnam, where they are known as xe ôm. Passengers can get a ride via mobile app or by hailing passing operators, or by finding drivers who gather at public places such as schools, markets, hospitals, and bus and train stations. Before the rise in popularity of ride-hailing apps, motorcycle taxi driving was a mostly informal economy, although some unions existed. Fare is verbally agreed upon before the trip based on distance. Some informal motorcycle taxi drivers still exist, as well as drivers working for regulated ride-hailing companies who would take on ad-hoc trips not booked through the app. Wearing a helmet on motorcycles is required by Vietnamese laws for both drivers and passengers, as such motorcycle taxi drivers would provide helmet for their customers. Go-Viet had a 35% market share among motorbike vehicle for hire companies in Ho Chi Minh City just six weeks after launching there on August 1, 2018, according to Go-Jek founder and chief executive Nadiem Makarim.	Technology	Motorized road transport	null
4048154	https://en.wikipedia.org/wiki/Si-o-se-pol	Si-o-se-pol	The Allahverdi Khan Bridge (), popularly known as Si-o-se-pol (), is the largest of the eleven historical bridges on the Zayanderud, the largest river of the Iranian Plateau, in Isfahan, Iran. The bridge was built in the early 17th century to serve as both a bridge and a dam. History Si-o-se-pol was built between 1599 and 1602, under the reign of Abbas the Great, the fifth shah of Safavid Iran. It was constructed under the supervision of Allahverdi Khan Undiladze, the commander-in-chief of the armies, who was of Georgian origin, and was also named after him. The bridge served particularly as a connection between the mansions of the elite, as well as a link to the city's vital Armenian neighborhood of New Julfa. In years of drought (2000–02 and 2013), the river was dammed upstream to provide water for Yazd province. Structure The bridge has a total length of and a total width of . It is a vaulted arch bridge consisting of two superimposed rows of 33 arches, from whence its popular name of Si-o-se-pol comes, and is made of stone. The longest span is about . The interior of Si-o-se-pol was originally decorated with paintings, which were often described by travelers as erotic. Gallery Transportation Chaharbagh Street Motahari Street Kamaloddin Esmaeil Street Chahar Bagh Bala Street Mellat Street Ayenekhaneh Street Enqelab Metro Station Si-o-se Pol Metro Station	Technology	Bridges	null
4050647	https://en.wikipedia.org/wiki/Pandanaceae	Pandanaceae	Pandanaceae is a family of flowering plants native to the tropics and subtropics of the Old World, from West Africa to the Pacific. It contains 982 known species in five genera, of which the type genus, Pandanus, is the most important, with species like Pandanus amaryllifolius and karuka (Pandanus julianettii) being important sources of food. The family likely originated during the Late Cretaceous. Characteristics Pandanaceae includes trees, shrubs, lianas, vines, epiphytes, and perennial herbs. Stems may be simple or bifurcately branched, and may have aerial prop roots. The stems bear prominent leaf scars. The leaves are very long and narrow, sheathing, simple, undivided, with parallel veins; the leaf margins and abaxial midribs are often prickly. The plants are dioecious. The inflorescences are terminally borne racemes, spikes or umbels, with subtended spathes, which may be brightly colored. The flowers are minute and lack perianths. Male flowers contain numerous stamens with free or fused filaments. Female flowers have a superior ovary, usually of many carpels in a ring, but may be reduced to a row of carpels or a single carpel. Fruits are berries or drupes, usually multiple. Pandanaceae includes five genera: Benstonea, Freycinetia, Martellidendron, Pandanus, and Sararanga. Benstonea (as subgenus "Acrostigma") and Martellidendron were formerly considered subgenera of Pandanus, but were recognized as distinct genera based on DNA sequencing. Uses Particular species of Pandanus are used to make mats (e.g. Central Africa) or in food products (e.g. leaves as flavoring, or fruit in Southeast Asia).	Biology and health sciences	Pandanales	Plants
7059576	https://en.wikipedia.org/wiki/Back-arc%20basin	Back-arc basin	A back-arc basin is a type of geologic basin, found at some convergent plate boundaries. Presently all back-arc basins are submarine features associated with island arcs and subduction zones, with many found in the western Pacific Ocean. Most of them result from tensional forces, caused by a process known as oceanic trench rollback, where a subduction zone moves towards the subducting plate. Back-arc basins were initially an unexpected phenomenon in plate tectonics, as convergent boundaries were expected to universally be zones of compression. However, in 1970, Dan Karig published a model of back-arc basins consistent with plate tectonics. Structural characteristics Back-arc basins are typically very long and relatively narrow, often thousands of kilometers long while only being a few hundred kilometers wide at most. For back-arc extension to form, a subduction zone is required, but not all subduction zones have a back-arc extension feature. Back-arc basins are found in areas where the subducting plate of oceanic crust is very old. The restricted width of back-arc basins is due to magmatic activity being reliant on water and induced mantle convection, limiting their formation to along subduction zones. Spreading rates vary from only a few centimeters per year (as in the Mariana Trough), to 15 cm/year in the Lau Basin. Spreading ridges within the basins erupt basalts that are similar to those erupted from the mid-ocean ridges; the main difference being back-arc basin basalts are often very rich in magmatic water (typically 1–1.5 weight % H2O), whereas mid-ocean ridge basalt magmas are very dry (typically <0.3 weight % H2O). The high water contents of back-arc basin basalt magmas is derived from water carried down the subduction zone and released into the overlying mantle wedge. Additional sources of water could be the eclogitization of amphiboles and micas in the subducting slab. Similar to mid-ocean ridges, back-arc basins have hydrothermal vents and associated chemosynthetic communities. Seafloor spreading Evidence of seafloor spreading has been seen in cores of the basin floor. The thickness of sediment that collected in the basin decreased toward the center of the basin, indicating a younger surface. The idea that thickness and age of sediment on the sea floor is related to the age of the oceanic crust was proposed by Harry Hess. Magnetic anomalies of the crust that had formed in back-arc basins deviated in form from the crust formed at mid-ocean ridges. In many areas the anomalies do not appear parallel, as well as the profiles of the magnetic anomalies in the basin lacking symmetry or a central anomaly as a traditional ocean basin does, indicating asymmetric seafloor spreading. This has prompted some to characterize the spreading in back-arc basins to be more diffused and less uniform than at mid-ocean ridges. The idea that back-arc basin spreading is inherently different from mid-ocean ridge spreading is controversial and has been debated through the years. Another argument put forward is that the process of seafloor spreading is the same in both cases, but the movement of seafloor spreading centers in the basin causes the asymmetry in the magnetic anomalies. This process can be seen in the Lau back-arc basin. Though the magnetic anomalies are more complex to decipher, the rocks sampled from back-arc basin spreading centers do not differ very much from those at mid-ocean ridges. In contrast, the volcanic rocks of the nearby island arc differ significantly from those in the basin. Back-arc basins are different from normal mid-ocean ridges because they are characterized by asymmetric seafloor spreading, but this is quite variable even within single basins. For example, in the central Mariana Trough, current spreading rates are 2–3 times greater on the western flank, whereas at the southern end of the Mariana Trough the position of the spreading center adjacent to the volcanic front suggests that overall crustal accretion has been nearly entirely asymmetric there. This situation is mirrored to the north where a large spreading asymmetry is also developed. Other back-arc basins such as the Lau Basin have undergone large rift jumps and propagation events (sudden changes in relative rift motion) that have transferred spreading centers from arc-distal to more arc-proximal positions. Conversely, study of recent spreading rates appear to be relatively symmetric with perhaps small rift jumps. The cause of asymmetric spreading in back-arc basins remains poorly understood. General ideas invoke asymmetries relative to the spreading axis in arc melt generation processes and heat flow, hydration gradients with distance from the slab, mantle wedge effects, and evolution from rifting to spreading. Formation and tectonics The extension of the crust behind volcanic arcs is believed to be caused by processes in association with subduction. As the subducting plate descends into the asthenosphere it sheds water, causing mantle melting, volcanism, and the formation of island arcs. Another result of this is a convection cell is formed. The rising magma and heat along with the outwards tension in the crust in contact with the convection cell cause a region of melt to form, resulting in a rift. This process drives the island arc toward the subduction zone and the rest of the plate away from the subduction zone. The backward motion of the subduction zone relative to the motion of the plate which is being subducted is called trench rollback (also known as hinge rollback or hinge retreat). As the subduction zone and its associated trench pull backward, the overriding plate is stretched, thinning the crust and forming a back-arc basin. In some cases, extension is triggered by the entrance of a buoyant feature in the subduction zone, which locally slows down subduction and induces the subducting plate to rotate adjacent to it. This rotation is associated with trench retreat and overriding plate extension. The age of the subducting crust needed to establish back-arc spreading has been found to be 55 million years old or older. This is why back-arc spreading centers appear concentrated in the western Pacific. The dip angle of the subducting slab may also be significant, as is shown to be greater than 30° in areas of back-arc spreading; this is most likely because as oceanic crust gets older it becomes denser, resulting in a steeper angle of descent. The thinning of the overriding plate from back-arc rifting can lead to the formation of new oceanic crust (i.e., back-arc spreading). As the lithosphere stretches, the asthenosphere below rises to shallow depths and partially melts as a result of adiabatic decompression melting. As this melt nears the surface, spreading begins. Sedimentation Sedimentation is strongly asymmetric, with most of the sediment supplied from the active volcanic arc which regresses in step with the rollback of the trench. From cores collected during the Deep Sea Drilling Project (DSDP) nine sediment types were found in the back-arc basins of the western Pacific. Debris flows of thick to medium bedded massive conglomerates account for 1.2% of sediments collected by the DSDP. The average size of the sediments in the conglomerates are pebble sized but can range from granules to cobbles. Accessory materials include limestone fragments, chert, shallow water fossils and sandstone clasts. Submarine fan systems of interbedded turbidite sandstone and mudstone made up 20% of the total thickness of sediment recovered by the DSDP. The fans can be divided into two sub-systems based on the differences in lithology, texture, sedimentary structures, and bedding style. These systems are inner and midfan subsystem and the outer fan subsystem. The inner and midfan system contains interbedded thin to medium bedded sandstones and mudstones. Structures that are found in these sandstones include load clasts, micro-faults, slump folds, convolute laminations, dewatering structures, graded bedding, and gradational tops of sandstone beds. Partial Bouma sequences can be found within the subsystem. The outer fan subsystem generally consists of finer sediments when compared to the inner and midfan system. Well sorted volcanoclastic sandstones, siltstones and mudstones are found in this system. Sedimentary structures found in this system include parallel laminae, micro-cross laminae, and graded bedding. Partial Bouma sequences can be identified in this subsystem. Pelagic clays containing iron-manganese micronodules, quartz, plagioclase, orthoclase, magnetite, volcanic glass, montmorillonite, illite, smectite, foraminiferal remains, diatoms, and sponge spicules made up the uppermost stratigraphic section at each site it was found. This sediment type consisted of 4.2% of the total thickness of sediment recovered by the DSDP. Biogenic pelagic silica sediments consist of radiolarian, diatomaceous, silicoflagellate oozes, and chert. It makes up 4.3% of the sediment thickness recovered. Biogenic pelagic carbonates is the most common sediment type recovered from the back-arc basins of the western Pacific. This sediment type made up 23.8% of the total thickness of sediment recovered by the DSDP. The pelagic carbonates consist of ooze, chalk, and limestone. Nanofossils and foraminifera make up the majority of the sediment. Resedimented carbonates made up 9.5% of the total thickness of sediment recovered by the DSDP. This sediment type had the same composition as the biogenic pelagic carbonated, but it had been reworked with well-developed sedimentary structures. Pyroclastics consisting of volcanic ash, tuff and a host of other constituents including nanofossils, pyrite, quartz, plant debris, and glass made up 9.5% of the sediment recovered. These volcanic sediments were sourced form the regional tectonic controlled volcanism and the nearby island arc sources. Locations Active back-arc basins are found in the Marianas, Kermadec-Tonga, South Scotia, Manus, North Fiji, and Tyrrhenian Sea regions, but most are found in the western Pacific. Not all subduction zones have back-arc basins; some, like the central Andes, are associated with rear-arc compression. There are a number of extinct or fossil back-arc basins, such as the Parece Vela-Shikoku Basin, Sea of Japan, and Kurile Basin. Compressional back-arc basins are found, for example, in the Pyrenees and the Swiss Alps. History of thought With the development of plate tectonic theory, geologists thought that convergent plate margins were zones of compression, thus zones of strong extension above subduction zones (back-arc basins) were not expected. The hypothesis that some convergent plate margins were actively spreading was developed by Dan Karig in 1970, while a graduate student at the Scripps Institution of Oceanography. This was the result of several marine geologic expeditions to the western Pacific.	Physical sciences	Tectonics	Earth science
6905166	https://en.wikipedia.org/wiki/Shimeji	Shimeji	Shimeji (Japanese: , or ) is a group of edible mushrooms native to East Asia, but also found in northern Europe. Hon-shimeji (Lyophyllum shimeji) is a mycorrhizal fungus and difficult to cultivate. Other species are saprotrophs, and buna-shimeji (Hypsizygus tessulatus) is now widely cultivated. Shimeji is rich in umami-tasting compounds such as guanylic acid, glutamic acid, and aspartic acid. Species Several species are sold as shimeji mushrooms. All are saprotrophic except Lyophyllum shimeji. Mycorrhizal Hon-shimeji (), Lyophyllum shimeji The cultivation methods have been patented by several groups, such as Takara Bio and Yamasa, and the cultivated hon-shimeji is available from several manufacturers in Japan. Saprotrophic Buna-shimeji (, lit. beech shimeji), Hypsizygus tessulatus, also known in English as the brown beech or brown clamshell mushroom. Hypsizygus marmoreus is a synonym of Hypsizygus tessulatus. Cultivation of Buna-shimeji was first patented by Takara Shuzo Co., Ltd. in 1972 as hon-shimeji and the production started in 1973 in Japan. Now, several breeds are widely cultivated and sold fresh in markets. Bunapi-shimeji (), known in English as the white beech or white clamshell mushroom. Bunapi was selected from UV-irradiated buna-shimeji ('hokuto #8' x 'hokuto #12') and the breed was registered as 'hokuto shiro #1' by Hokuto Corporation. Hatake-shimeji (), Lyophyllum decastes. Shirotamogidake (), Hypsizygus ulmarius. These two species had been also sold as hon-shimeji. Velvet pioppino (alias velvet pioppini, black poplar mushroom, Chinese: /), Agrocybe aegerita. Shimeji health benefits Shimeji mushrooms contain minerals like potassium and phosphorus, magnesium, zinc, and copper. Shimeji mushrooms lower the cholesterol level of the body. This mushroom is rich in glycoprotein (HM-3A), marmorin, beta-(1-3)-glucan, hypsiziprenol, and hypsin therefore is a potential natural anticancer agent. Shimeji mushrooms contain angiotensin I-converting enzyme (ACE) inhibitor which is an oligopeptide that may be helpful in lowering blood pressure and reducing the risk of stroke in persons having hypertension. Also rich in polysaccharides, phenolic compounds, and flavonoids. Therefore, inhibits inflammatory cytokines and oxidative stress and protects from lung failure. These compounds also help in reducing oxidative stress-mediated disease through radical scavenging activity hence these mushrooms are antioxidants also. Culinary Use Shimeji should always be cooked: it is not a good mushroom to serve raw due to a somewhat bitter taste, but the bitterness disappears completely upon cooking. The cooked mushroom has a pleasant, firm, slightly crunchy texture and a slightly nutty flavor. Cooking also makes this mushroom easier to digest. It works well in stir-fried foods like stir-fried vegetables, as well as with wild game or seafood. Also, it can be used in soups, stews, and in sauces. When cooked alone, Shimeji mushrooms can be sautéed whole, including the stem or stalk (only the very end cut off), using a higher temperature or they can be slow roasted at a low temperature with a small amount of butter or cooking oil. Shimeji is used in soups, nabe and takikomi gohan.	Biology and health sciences	Edible fungi	Plants
9174535	https://en.wikipedia.org/wiki/Aquaculture%20of%20tilapia	Aquaculture of tilapia	Tilapia has become the third most important fish in aquaculture after carp and salmon; worldwide production exceeded in 2002 and increases annually. Because of their high protein content, large size, rapid growth (6 to 7 months to grow to harvest size), and palatability, a number of coptodonine and oreochromine cichlids—specifically, various species of Coptodon, Oreochromis, and Sarotherodon—are the focus of major aquaculture efforts. Tilapia fisheries originated in Africa and the Levant. The accidental and deliberate introductions of tilapia into South and Southeast Asian freshwater lakes have inspired outdoor aquaculture projects in various countries with tropical climates, including Honduras, Papua New Guinea, the Philippines, and Indonesia. Tilapia farm projects in these countries have the highest potential to be "green" or environmentally friendly. In temperate zone localities, tilapia farmers typically need a costly energy source to maintain a tropical temperature range in their tanks. One relatively sustainable solution involves warming the tank water using waste heat from factories and power stations. Tilapiines are among the easiest and most profitable fish to farm due to their omnivorous diet, mode of reproduction (the fry do not pass through a planktonic phase), tolerance of high stocking density, and rapid growth. In some regions the fish can be raised in rice fields at planting time and grow to edible size () when the rice is ready for harvest. Unlike salmon, which rely on high-protein feeds based on fish or meat, commercially important tilapiine species eat a vegetable or cereal-based diet. Tilapia raised in inland tanks or channels are considered safe for the environment, since their waste and disease is contained and not spread to the wild. However, tilapiines have acquired notoriety as being among the most serious invasive species in many subtropical and tropical parts of the world. For example, blue tilapia (Oreochromis aureus) (itself commonly confused with another species often used in aquaculture, the Nile tilapia, O. niloticus), Mozambique tilapia (O. mossambicus), blackchin tilapia (Sarotherodon melanotheron), spotted tilapia (Pelmatolapia mariae), and redbelly tilapia (Coptodon zillii) have all become established in the southern United States, particularly in Florida and Texas. Commercially grown tilapia are almost exclusively male. Being prolific breeders, female tilapia in the ponds or tanks will result in large populations of small fish. Whole tilapia can be processed into skinless, boneless (PBO) fillets: the yield is from 30% to 37%, depending on fillet size and final trim. Commercial breeding of Nile tilapia Although farming of Tilapia has been going on for thousands of years, the breeding of Tilapia did not start until recently. The first breeding program started in 1988 in a collaboration between the international center for living aquatic resources (ICLARM or WorldFish) and Akvaforsk. The name of the project was GIFT, meaning genetically improved farmed tilapia. Four wild strains from Africa were crossed with four farmed strains from the Philippines. This strain is currently farmed in more than 87 countries in Asia, Africa and Latin America. The GIFT strain is used in two selection programs, one of them being GenoMar, a subsidiary EW Group. In the past the absolute and only important trait when breeding tilapia was growth, being the only criteria for selection. Today more traits have been added into the selection criteria, like growth, fillet yield, robustness and specific disease resistance. Robustness is one of those traits that is becoming more important since it is the biggest problem with mortality on farms today. GenoMar has successfully had a growth increase of 7% per generation while fillet yield only improves with 0,3% per generation. The explanation for this is its low heritability together with the fact that the trait cannot be measured on live animals and therefore information of fillet yield is given from relatives instead. Breeding of Tilapia is done with the help of a pyramid scheme with multiplying generations. The goal with this is that a few high merit individuals can be passed down into billions of production fish at the farms. The generation interval today is down to only 6-9 months meaning that there can be more than one generation per year. Mass selection and pedigree-based selection are the most used methods today for genetic improvements of tilapias. The breeding program GenoMar has used marker-assisted selection since 2004 using microsatellites when doing parentage assignment has been done on Tilapia. Since 2019 genomic selection using single nucleotide polymorphic (SNP) has been used more widely. The latest genome assembly is from 3 years ago, and can be seen from the name of the assembly that NMBU was also highly involved in the effort along with the University of Maryland. Nile tilapia have 22 pairs of chromosomes. 23 Linkage groups because of the sex chromosomes. About 1 billion base pairs in length, 3,010 contigs made 2,460 scaffolds which were placed at a chromosome level as the karyotype was already known. The reference genome also has the non-nuclear mitochondrial genome. With help from bioinformatics, the estimated number of genes is around 30 thousand. It can be browsed in ensembl or NCBI for example. The essential technological foundation for genomic selection is not obtaining the pedigree but the genotypes of the animals. This is currently done by SNP chips, oligonucleotide arrays. They have to be specifically designed for a species and commonly that is done based after the whole genome sequence has been obtained. Only some parts of the WGS are of interest, namely those that exhibit variation. Dense markers are considered good enough to capture the gene content because they are in linkage disequilibrium with the genes that influence the phenotype, or as they are called QTL. There are three major SNP chips for Tilapia, which were announced in 2018, 2020 and then 2020 again. In RNA Seq, the RNA from cells of a tissue is extracted and they are sequenced in order to know what genes are being expressed and at what intensities. First, the extracted RNA is converted into cDNA and that cDNA library is sequenced using the same machines used in whole genome sequencing. The bioinformatics pipeline afterwards is different from WGS, though. Since the segments are shorter, alignment is not as difficult. Also, one is interested in the amount of transcription that is happening. In any case, RNA Seq provides valuable insight into the biology of a group of cells and the entire organism. In the beginning, breeding programs focused mainly on growth traits. Nowadays, more traits are included in the breeding goal. The tendency is to have even more traits in the future. For example, disease resistance, reproductive traits, robustness, lower emissions, feed conversion ratio (sustainability traits). New technologies for high-throughput phenotyping, as in the concept of precision farming, mean that many novel traits might be included as well. There is already a tendency to have mergers and acquisitions, with a few companies buying smaller ones, just like it happened in the poultry industry. Also like the poultry industry, the exploitation of heterosis is a possibility that could be established in Tilapia breeding. This would mean more protection for the companies and attract more corporate enterprises. Because crossbreeding provides a biological lock mechanism. Gene editing is coming. Especially CRISPR-Cas9 can already be implemented and holds a lot of promise. We also expect for the production systems to become ever more intensive. This means fewer ponds, more cages and more RAS. Overall, the average fish density is going to increase. Nutritional value Tilapia from aquaculture contain especially high ratios of omega-6 to omega-3 fatty acids. Around the world Apart from the very few species found in the Western Asia, such as the Middle Eastern mango tilapia, there are no tilapia cichlids native to Asia. However, species originally from Africa have been widely introduced and have become economically important as food fish in many countries. China, the Philippines, Taiwan, Indonesia, and Thailand are the leading suppliers, and these countries altogether produced about of fish in 2001, constituting about 76% of the total aquaculture production of tilapia worldwide. Other countries India The FAO has not recorded any production of farmed tilapia by India. Rajiv Gandhi Centre for Aquaculture (RGCA), the R&D arm of Marine Products Export Development Authority, has established a facility in Vijayawada to produce mono-sex tilapia in two strains. This project involves the establishment of a satellite nucleus for the GIFT strain of tilapia in India, the design and conduct of a genetic improvement program for this strain, the development of dissemination strategies, and the enhancement of local capacity in the areas of selective breeding and genetics. The development and dissemination of a high yielding tilapia strain possessing desirable production characteristics is expected to bring about notable economic benefits for the country. Farming of Tilapia is not permitted in the country on commercial basis. The Rajiv Gandhi Center for Aquaculture (RGCA) has expressed interest in obtaining the Genetically Improved Farmed Tilapia (GIFT strain) for aquaculture development in the country. The GIFT tilapia strain, selectively bred in Malaysia and the Philippines, has achieved an improvement of more than 10 per cent per generation in growth rate and has been widely distributed to several Asian countries and to Latin America (Brazil). However, rather than passively importing the improved genetic stock, the Center is interested in running a formal breeding program (fully pedigreed population) similar to the one that has been carried out for the GIFT strain in Malaysia. The aim is to produce fast-growing high yielding tilapia strains adapted to a wide range of local farming environments that can be grown at as low a cost as possible. The project involves several steps. The first is the establishment of a new nucleus of the GIFT strain at the RGCA and the design of a formal breeding program to further improve its genetic performance within the local environment. This will involve enhancing the capacity of local personnel in selective breeding, genetic improvement, statistical analysis and hatchery management through specialized training courses. Once a high performing tilapia strain (or strains) has been developed, the establishment of satellite hatcheries will increase the availability and decrease the costs of seed stock. These public and private hatcheries will act as multipliers for the superior genetics developed at RGCA and the sites for dissemination of quality broodstock to fish farmers. Although the ultimate target groups of this project are fish farmers and small householders, a wider range of beneficiaries is expected, including commercial producers, scientists and the end consumers. The RGCA will gain experience and knowledge on the development of genetic improvement programs for economically important traits and other aspects of modern quantitative genetics. This experience and the development of a standard selective breeding protocol will allow for genetic improvement programs for other aquaculture species that are commonly cultured in India. Hatchery managers, producers and farmers will also improve their capacity to implement on-farm selective breeding programs. In the longer term the project is also expected to contribute to the development of a complete chain of production. This will require initial capital support for farmers, identification of alternative cheap plant-based feed, and diagnosis of diseases in hatcheries, as well as strategies for early growth management. Improvement in harvest technologies, including storage of product and transport facilities, is likely to improve as a consequence of this project. Malawi In 2010 Malawi produced 2,997 tonnes of farmed tilapia. A few species of Oreochromis tilapia, popular known as 'chambo', are the most popular fish in Malawi. They are endemic to bodies of water in Malawi like Lake Malawi, Lake Malombe and the Shire River. Due to over fishing, the fish however is now on the threatened species list. Malawi has fish farms that are dedicated to farming tilapia.	Technology	Aquaculture	null
10674423	https://en.wikipedia.org/wiki/Taobao	Taobao	Taobao is a Chinese online shopping platform. It is headquartered in Hangzhou and is owned by Alibaba. According to Alexa rank, it was the eighth most-visited website globally in 2021. Taobao.com was registered on April 21, 2003 by Alibaba Cloud Computing (Beijing) Co., Ltd. Taobao Marketplace facilitates consumer-to-consumer retail by providing a platform for small businesses and individual entrepreneurs to open online stores that mainly cater to consumers in Chinese-speaking regions (Mainland China, Hong Kong, Macau and Taiwan) and abroad, which is made payable by online accounts. Its stores usually offer an express delivery service. Sellers are able to post goods for sale either through a fixed price or an auction. Auctions make up a small percentage of transactions, whereas the majority of the products are new merchandise sold at fixed prices. Taobao users usually read feedback and compare items from multiple shops. Taobao's popular payment platform is Alibaba's Alipay As of at least 2024, it is the world's most popular shopping hub as measured by gross merchandise value. History Before its launch of Taobao, Alibaba had focused on online business-to-business wholesale sales. In 2003, eBay acquired Eachnet, China's online auction leader at the time, for US$180 million. It became a major contender in the Chinese consumer e-commerce market. Responding to eBay's moves Alibaba launched Taobao as a rival consumer-to-consumer platform. To counter eBay's expansion, Taobao offered free listings to sellers. It introduced instant messaging for facilitating buyer-seller communication and an escrow-based payment tool: Alipay. Taobao's focus on institutional trust building mechanisms like escrowing payments became a major reason for its success in the market for eBay, despite eBay's first-mover advantage. Taobao became mainland China's market leader within two years. Its market share grew from 8% to 59% between 2003 and 2005, while eBay China dropped from 79% to 36%. eBay shut down its Chinese site in 2006. In 2008, Taobao established a platform rule providing that customers had the right to return clothes sold on the platform within seven days of receipt without cause, and subsequently expanded the rule to cover other commodities. This rule became an influential standard in Chinese e-commerce and in 2014 was made an industry standard through the State Administration for Industry and Commerce's Administrative Measures for Online Trading. In October 2010, Taobao beta-launched eTao as an independent search engine for online shopping to provide and merchant information from a number of major consumer e-commerce websites in China. Online shoppers would be able to use the site to compare prices across sellers. According to the Alibaba Group web site, eTao offers products from Amazon China, Dangdang, Gome, Yihaodian, Nike China and Vancl, as well as Taobao and Tmall. In June 2011, Jack Ma, executive chairman and former chief executive officer of Alibaba Group, announced that Taobao would split into three different companies: Taobao Marketplace (a consumer-to-consumer platform), Tmall.com (a business-to consumer platform, then called Taobao Mall), and eTao (a search engine for online shopping). The move was said to be necessary for Taobao to “meet competitive threats that emerged in the past two years during which the Internet and e-commerce landscape has changed dramatically.” In 2012 Taobao began to accept international Visa and MasterCard credit and debit cards. On 29 April 2013, Alibaba announced an investment of US$586 million in Sina Weibo. According to Reuters, the deal “should drive more web traffic to Alibaba's Taobao Marketplace”. On August 1, 2013, Alibaba launched Weibo for Taobao, which allows users to link Sina Weibo accounts with Taobao accounts. In addition to hosting individuals and businesses, Taobao includes online stores for courts, customs offices, state-owned banks, and asset management companies selling distressed assets. By early 2014, more than 500 local Chinese courts had established Taobao store fronts to sell seized property, including property which had been confiscated as part of the anti-corruption campaign under Xi Jinping. Features Shop feedback A good way to investigate a Taobao shop is by clicking the shop's rating icon. For Tmall.com shops, people click the stars to view their ratings. Taobao users usually read feedback and compare items from multiple shops. Feedback can be genuine or artificial, requiring users to make their own judgments. Feedback can be posted by competitors. Every trade deal includes a section of customer feedback. Shop owners often put effort to maximize positive comments. Negotiations may happen between sellers and consumers over their satisfaction ratings. Taobao uses search tools and other functions to understand user demands. Customers are asked to complete surveys that ask: Obstacles encountered Advice and reviews Use frequency Importance of that particular product Satisfaction levels about specific aspects, including visual and typographic layout, procedure and instructions Overall satisfaction Other features Taobao Marketplace offers various features and services to create a better user experience for online shoppers and retailers. In January 2010, it launched the Taobao app, created by independent developers through the Taobao Open Platform, to be downloaded by consumers in Taobao App Store. In March 2010, it introduced the Taobao Data Cube platform, which gives small businesses access to its aggregate consumer transactions data for insight into industry trends. In June 2010, it partnered with Wasu Media Internet Limited to launch Taohua, a digital entertainment products platform, and interactive digital television shopping, that are operated by a joint-venture formed by the two companies. Weitao Weitao is a private shopping assistant/blog for Taobao/Tmall customers. It is a micro-blogging feature for brands and merchants on its e-commerce sites Taobao and Tmall. Taobao General Code Taobao has developed an extensive set of rules and a constitution which it terms the "General Code". The General Code consists of 6 chapters and 31 provisions proscribing the basic access requirements and obligations of users, seller obligations, and the platform's conflict resolution mechanism. When Taobao proposes a rule change, all buyers and sellers with a sufficiently high Zhima credit score can vote and express their opinions on the rule. Taobao also invites professionals and academics with relevant expertise to public evaluate proposed rule changes in order to inform voters. Public jury Alibaba established its public jury (pan.taobao.com) dispute resolution system in 2012. By 2020, it was a widely-used dispute resolution tool on the main Taobao platform and was also frequently used on Xianyu (Taobao's used goods platform). Through the public jury process, Taobao randomly selects panels of 13 jurors (termed "public assessors" on Taobao) from a pool of 4 million volunteers. Candidates must have been on Taobao for at least a year and have sufficiently high Sesame Credit ratings. Volunteers earn experience points that can earn virtual titles and which can be "spent" for Taobao to make a donation to charity. All participants in the public jury process are anonymized and no communication between or among the disputing parties and the jurors are permitted. Jurors review the case and vote within 48 hours. The party with the most votes wins. A party who is unsatisfied with the jury outcome can request further review by Taobao employees. Taobao can implement jury decisions through means including freezing payments, taking money from a sellers' store deposit, lowering user ratings, or removing a party from the platform. Most public jury cases involve buyer-seller disputes. In some instances, Taobao has used large juries of 800 to 1000 jurors to decide issues relating to platform governance. For example, Taobao used a large jury to decide whether to allow a baby bottle manufacturer, Betta, to remain on the platform. Betta was a legal but copycat style product of the popular Japanese brand Doctor Betta, which was also sold on Taobao. Jurors voted to remove the copycat product from the platform. Golden Cudgel In 2019, Taobao launched a dispute resolution called Golden Cudgel, named after the monkey king's magical weapon in the 16th century novel Journey to the West. This mechanism allows sellers who have passed qualifying exams to remove a limited number of malicious reviews on a daily basis without prior permission from the platform. Sellers must submit evidence to establish the dishonest nature of each removed review to enable review by the platform. Sellers' ability to use the Golden Cudgel mechanism is revoked if they are found by the platform to have repeatedly removed reviews improperly. Services Alipay Launched in 2004, Alipay () is an escrow-based online payment platform. It is the preferred payment solution for Taobao Marketplace. It is the most widely used third-party online payment solution in China. To ensure safe transactions, Alipay uses an escrow system through which payment is only released to the seller once the buyer has received goods in satisfactory condition. According to the Alibaba Group website, Alipay partners with multiple financial institutions as well as Visa and MasterCard to facilitate payments in China and abroad. At the time it was implemented, the mandatory escrow feature of Alipay was a major institutional innovation for e-commerce platforms. This was a major reason Taobao was able to outcompete eBay/EachNet in the Chinese market. Alipay systems are separate for different groups of users. For instance, Alipay users may send and receive funds if they have an account that have a credit card issued in China, whereas users with other cards may only use Alipay to pay for goods or services from Taobao. This has proven problematic for international users, as they are unable to receive refunds not issued via the Taobao system. AliExpress AliExpress was created in April 2010 as an international retailing website. People who live overseas can use the service to purchase items from Chinese manufacturers online. AliWangWang (TradeManager) Taobao Marketplace allows buyers and sellers to communicate prior to the purchase through its embedded proprietary instant chat program, named AliWangWang (). Baopals In February 2016, 3 expats living in Shanghai launched Baopals, a shopping platform that translates Taobao and Tmall into English so that foreigners living in China can access of its products and services. Over 2.4 million items have been sold on the platform. Happy Taobao In Dec. 2009, Taobao, together with Hunan TV, set up Happy Taobao, Inc for television shopping. Hunan TV launched an entertainment series called "Happy Taobao", while Taobao Marketplace created channels and independent websites. Taobao Live In 2018, Alibaba launched the streaming service named Alibaba Live. This service was created with the goal of allowing online retailers to market their products utilizing social shopping. This has seen significant growth in popularity and success, with the 84 stores using this service reporting $7.4 million in 2020 sales. Taobao stated that they predict live-streaming on their platform will generate over 500 billion sales transactions. Taobao Live now has over 10,000 weblebrities promoting a wide range of items including as cosmetics, apparel, cuisine, and numerous electrical gadgets. Taobao Live's daily sales have already surpassed $3 billion. Alibaba promotes a new style of live streaming, called cūnbō (村播), that features rural sellers. Taobao has given them their own category in the app, with the purpose of making it easier for these rural sellers to find customers and followers on the platform. Singles' Day Singles' Day (also known as the Double eleven shopping carnival, as in 11/11) is the largest Chinese online shopping day. It takes place on 11 November each year. It takes the advantages of the Chinese singles day that was created by Chinese university students to celebrate their bachelordom. After the event was launched, it obtained widespread attention, attracting other e-commerce companies to imitate this model. Singles day grew rapidly since its introduction in 2009. 2009 sales reached RMB50 million (£5.68 million): Sales grew rapidly thereafter: 2009: RMB50 million (£5.68 million) 2010: RMB900 million (£102 million) 2011: RMB3.4 billion (£386 million) 2012: RMB19.1 billion (£2.17 billion) 2013: RMB35 billion (£3.97 billion) 2015: RMB91.2 billion 2016: RMB120.7 billion 2017: RMB168.2 billion 2018: RMB213.5 billion 2019: RMB268.4 billion In 2016, Alibaba introduced the T-mall double eleven party, inviting celebrities who took part in a Victoria's Secret show. At the 2017 party, Jack Ma launched his film Gong Shou Dao (Defend the Homeland with Kungfu). Because of the huge trading volume and income in Singles Day, Taobao launched another promotional activity on December 12 (12/12), drove record trading every year thereafter. Technology Anti-fraud Taobao has a five layer system for fraud detection. "Account check" is the first layer: during this stage, automatic processes examine whether the account at issue has demonstrated suspicious activity. Fraud cases deemed obvious at this layer are declined. Taobao's system can also require further information and submit the transaction to the next three layers of automatic review: device check, activity check, and risk strategy. Each check sends fraud cases deemed obvious to an automatic decision and refers potential no-obvious fraud to the next level. These automated anti-fraud checks use big data models, including analysis of user behavioral data, network data, delivery details, and IP addresses. The final level is manual review by Taobao employees. Taobao shares information with police and local courts to assist in locating sellers alleged to be selling counterfeit products. For example, in 2014, collaboration with law enforcement resulted in 1,000 counterfeit cases, 400 arrests, and the shutdown of 200 related physical stores, warehouses, and factories. Markets Taobao for Southeast Asia In September 2013, Taobao launched its Southeast Asian site. A translation feature is available for major languages in Southeast Asia. In September 2024, Taobao launched an English version of its app tailored specifically for users in Singapore. Controversies Taobao has sometimes been the subject of in-person and online seller protests following major changes to its rules. The largest seller protest was the 2011 "October Rising". With the goal of reducing counterfeits and substandard products, Taobao had increased the Taobao Mall membership fees for sellers and their required cash deposits. The rule changes were made without warning. Approximately 50,000 sellers formed the "anti-Taobao alliance" for digital protest actions and in-person protest at Alibaba's headquarters. The Chinese government mediated the dispute, resulting in Taobao revising its seller fees and providing 1.8 billion RMB in support for small businesses using the platform. In August 2017, the company removed controversial vendors offering personalised messages featuring African children over concerns of child exploitation. Some Chinese Taobao vendors claimed that their promotional videos featuring African children were "charity activity" in which most of the profits goes to the children. However the situation proved more complicated after a photographer contacted by the Beijing Youth Daily said "the children only received snacks or a few dollars as reward", indicating that there was legitimate child exploitation. In 2019, Taobao removed all items related to the Houston Rockets in response to the organization's general manager Daryl Morey posting a tweet about Hong Kong. In October 2020 amid rising geopolitical tensions between Taipei and Beijing, Taobao announced that it would exit the Taiwanese market after the Taiwanese government ordered the company to re-register as being backed by China or to leave the island if they don't. The mobile app of Taobao was banned in India (along with other Chinese apps) on 2 September 2020 by the government, the move came amid the 2020 China-India skirmish. In 2022, the Office of the United States Trade Representative named Taobao on its list of Notorious Markets for Counterfeiting and Piracy. Metrics Taobao Marketplace had more than 5 million registered users as of June 2013 and hosted more than 80 million product listings. It facilitated approximately RMB 200 billion in gross merchandise volume in 2009. In September 2013, Taobao ranked 12th overall in Alexa's internet rankings. With over 1 billion product listings as of 2016, the combined transaction volume of Taobao Marketplace and Tmall.com reached 3 trillion yuan in 2017. As of 2020, Taobao hosted ten million online stores, 726 million active buyers, and a gross merchandise value of $945 billion. As of 2021, Taobao was the 8th most visited website in the world and the 5th most visited website in China. As of at least 2024, it was the world's most popular online shopping platform as measured by gross merchandise value. Taobao villages Taobao villages are rural Chinese villages where the local economy has developed to focus extensively on Taobao. Alibaba's research division defines Taobao villages as those in which (1) businesses are located in an administrative village in a rural area, (2) the village's annual e-commerce revenues exceed RMB 10 million, and (3) the village has either an excess of 100 active online shops or active online shops account for more than 10% of village households.	Technology	E-commerce	null
10686210	https://en.wikipedia.org/wiki/Jupiter%20mass	Jupiter mass	The Jupiter mass, also called Jovian mass, is the unit of mass equal to the total mass of the planet Jupiter. This value may refer to the mass of the planet alone, or the mass of the entire Jovian system to include the moons of Jupiter. Jupiter is by far the most massive planet in the Solar System. It is approximately 2.5 times as massive as all of the other planets in the Solar System combined. Jupiter mass is a common unit of mass in astronomy that is used to indicate the masses of other similarly-sized objects, including the outer planets, extrasolar planets, and brown dwarfs, as this unit provides a convenient scale for comparison. Current best estimates The current best known value for the mass of Jupiter can be expressed as : which is about as massive as the Sun (is about ): Jupiter is 318 times as massive as Earth: Context and implications Jupiter's mass is 2.5 times that of all the other planets in the Solar System combined—this is so massive that its barycenter with the Sun lies beyond the Sun's surface at 1.068 solar radii from the Sun's center. Because the mass of Jupiter is so large compared to the other objects in the Solar System, the effects of its gravity must be included when calculating satellite trajectories and the precise orbits of other bodies in the Solar System, including the Moon and even Pluto. Theoretical models indicate that if Jupiter had much more mass than it does at present, its atmosphere would collapse, and the planet would shrink. For small changes in mass, the radius would not change appreciably, but above about (1.6 Jupiter masses) the interior would become so much more compressed under the increased pressure that its volume would decrease despite the increasing amount of matter. As a result, Jupiter is thought to have about as large a diameter as a planet of its composition and evolutionary history can achieve. The process of further shrinkage with increasing mass would continue until appreciable stellar ignition was achieved, as in high-mass brown dwarfs having around 50 Jupiter masses. Jupiter would need to be about 80 times as massive to fuse hydrogen and become a star. Gravitational constant The mass of Jupiter is derived from the measured value called the Jovian mass parameter, which is denoted with GMJ. The mass of Jupiter is calculated by dividing GMJ by the constant G. For celestial bodies such as Jupiter, Earth and the Sun, the value of the GM product is known to many orders of magnitude more precisely than either factor independently. The limited precision available for G limits the uncertainty of the derived mass. For this reason, astronomers often prefer to refer to the gravitational parameter, rather than the explicit mass. The GM products are used when computing the ratio of Jupiter mass relative to other objects. In 2015, the International Astronomical Union defined the nominal Jovian mass parameter to remain constant regardless of subsequent improvements in measurement precision of . This constant is defined as exactly If the explicit mass of Jupiter is needed in SI units, it can be calculated by dividing GM by G, where G is the gravitational constant. Mass composition The majority of Jupiter's mass is hydrogen and helium. These two elements make up more than 87% of the total mass of Jupiter. The total mass of heavy elements other than hydrogen and helium in the planet is between 11 and . The bulk of the hydrogen on Jupiter is solid hydrogen. Evidence suggests that Jupiter contains a central dense core. If so, the mass of the core is predicted to be no larger than about . The exact mass of the core is uncertain due to the relatively poor knowledge of the behavior of solid hydrogen at very high pressures. Relative mass	Physical sciences	Mass and weight	Basics and measurement
8556497	https://en.wikipedia.org/wiki/Contraposition	Contraposition	In logic and mathematics, contraposition, or transposition, refers to the inference of going from a conditional statement into its logically equivalent contrapositive, and an associated proof method known as . The contrapositive of a statement has its antecedent and consequent inverted and flipped. Conditional statement . In formulas: the contrapositive of is . If P, Then Q. — If not Q, Then not P. "If it is raining, then I wear my coat" — "If I don't wear my coat, then it isn't raining." The law of contraposition says that a conditional statement is true if, and only if, its contrapositive is true. Contraposition () can be compared with three other operations: Inversion (the inverse), "If it is not raining, then I don't wear my coat." Unlike the contrapositive, the inverse's truth value is not at all dependent on whether or not the original proposition was true, as evidenced here. Conversion (the converse), "If I wear my coat, then it is raining." The converse is actually the contrapositive of the inverse, and so always has the same truth value as the inverse (which as stated earlier does not always share the same truth value as that of the original proposition). Negation (the logical complement), "It is not the case that if it is raining then I wear my coat.", or equivalently, "Sometimes, when it is raining, I don't wear my coat. " If the negation is true, then the original proposition (and by extension the contrapositive) is false. Note that if is true and one is given that is false (i.e., ), then it can logically be concluded that must be also false (i.e., ). This is often called the law of contrapositive, or the modus tollens rule of inference. Intuitive explanation In the Euler diagram shown, if something is in A, it must be in B as well. So we can interpret "all of A is in B" as: It is also clear that anything that is not within B (the blue region) cannot be within A, either. This statement, which can be expressed as: is the contrapositive of the above statement. Therefore, one can say that In practice, this equivalence can be used to make proving a statement easier. For example, if one wishes to prove that every girl in the United States (A) has brown hair (B), one can either try to directly prove by checking that all girls in the United States do indeed have brown hair, or try to prove by checking that all girls without brown hair are indeed all outside the US. In particular, if one were to find at least one girl without brown hair within the US, then one would have disproved , and equivalently . In general, for any statement where A implies B, not B always implies not A. As a result, proving or disproving either one of these statements automatically proves or disproves the other, as they are logically equivalent to each other. Formal definition A proposition Q is implicated by a proposition P when the following relationship holds: This states that, "if , then ", or, "if Socrates is a man, then Socrates is human." In a conditional such as this, is the antecedent, and is the consequent. One statement is the contrapositive of the other only when its antecedent is the negated consequent of the other, and vice versa. Thus a contrapositive generally takes the form of: That is, "If not-, then not-", or, more clearly, "If is not the case, then P is not the case." Using our example, this is rendered as "If Socrates is not human, then Socrates is not a man." This statement is said to be contraposed to the original and is logically equivalent to it. Due to their logical equivalence, stating one effectively states the other; when one is true, the other is also true, and when one is false, the other is also false. Strictly speaking, a contraposition can only exist in two simple conditionals. However, a contraposition may also exist in two complex, universal conditionals, if they are similar. Thus, , or "All s are s," is contraposed to , or "All non-s are non-s." Sequent notation The transposition rule may be expressed as a sequent: where is a metalogical symbol meaning that is a syntactic consequence of in some logical system; or as a rule of inference: where the rule is that wherever an instance of "" appears on a line of a proof, it can be replaced with ""; or as the statement of a truth-functional tautology or theorem of propositional logic. The principle was stated as a theorem of propositional logic by Russell and Whitehead in Principia Mathematica as where and are propositions expressed in some formal system. Proofs Simple proof by definition of a conditional In first-order logic, the conditional is defined as: which can be made equivalent to its contrapositive, as follows: Simple proof by contradiction Let: It is given that, if A is true, then B is true, and it is also given that B is not true. We can then show that A must not be true by contradiction. For if A were true, then B would have to also be true (by Modus Ponens). However, it is given that B is not true, so we have a contradiction. Therefore, A is not true (assuming that we are dealing with bivalent statements that are either true or false): We can apply the same process the other way round, starting with the assumptions that: Here, we also know that B is either true or not true. If B is not true, then A is also not true. However, it is given that A is true, so the assumption that B is not true leads to a contradiction, which means that it is not the case that B is not true. Therefore, B must be true: Combining the two proved statements together, we obtain the sought-after logical equivalence between a conditional and its contrapositive: More rigorous proof of the equivalence of contrapositives Logical equivalence between two propositions means that they are true together or false together. To prove that contrapositives are logically equivalent, we need to understand when material implication is true or false. This is only false when is true and is false. Therefore, we can reduce this proposition to the statement "False when and not-" (i.e. "True when it is not the case that and not-"): The elements of a conjunction can be reversed with no effect (by commutativity): We define as equal to "", and as equal to (from this, is equal to , which is equal to just ): This reads "It is not the case that (R is true and S is false)", which is the definition of a material conditional. We can then make this substitution: By reverting R and S back into and , we then obtain the desired contrapositive: In classical propositional calculus system In Hilbert-style deductive systems for propositional logic, only one side of the transposition is taken as an axiom, and the other is a theorem. We describe a proof of this theorem in the system of three axioms proposed by Jan Łukasiewicz: A1. A2. A3. (A3) already gives one of the directions of the transposition. The other side, , is proven below, using the following lemmas proven here: (DN1) - Double negation (one direction) (DN2) - Double negation (another direction) (HS1) - one form of Hypothetical syllogism (HS2) - another form of Hypothetical syllogism. We also use the method of the hypothetical syllogism metatheorem as a shorthand for several proof steps. The proof is as follows: (instance of the (DN2)) (instance of the (HS1) (from (1) and (2) by modus ponens) (instance of the (DN1)) (instance of the (HS2)) (from (4) and (5) by modus ponens) (from (3) and (6) using the hypothetical syllogism metatheorem) (instance of (A3)) (from (7) and (8) using the hypothetical syllogism metatheorem) Comparisons Examples Take the statement "All red objects have color." This can be equivalently expressed as "If an object is red, then it has color." The contrapositive is "If an object does not have color, then it is not red." This follows logically from our initial statement and, like it, it is evidently true. The inverse is "If an object is not red, then it does not have color." An object which is blue is not red, and still has color. Therefore, in this case the inverse is false. The converse is "If an object has color, then it is red." Objects can have other colors, so the converse of our statement is false. The negation is "There exists a red object that does not have color." This statement is false because the initial statement which it negates is true. In other words, the contrapositive is logically equivalent to a given conditional statement, though not sufficient for a biconditional. Similarly, take the statement "All quadrilaterals have four sides," or equivalently expressed "If a polygon is a quadrilateral, then it has four sides." The contrapositive is "If a polygon does not have four sides, then it is not a quadrilateral." This follows logically, and as a rule, contrapositives share the truth value of their conditional. The inverse is "If a polygon is not a quadrilateral, then it does not have four sides." In this case, unlike the last example, the inverse of the statement is true. The converse is "If a polygon has four sides, then it is a quadrilateral." Again, in this case, unlike the last example, the converse of the statement is true. The negation is "There is at least one quadrilateral that does not have four sides." This statement is clearly false. Since the statement and the converse are both true, it is called a biconditional, and can be expressed as "A polygon is a quadrilateral if, and only if, it has four sides." (The phrase if and only if is sometimes abbreviated as iff.) That is, having four sides is both necessary to be a quadrilateral, and alone sufficient to deem it a quadrilateral. Truth If a statement is true, then its contrapositive is true (and vice versa). If a statement is false, then its contrapositive is false (and vice versa). If a statement's inverse is true, then its converse is true (and vice versa). If a statement's inverse is false, then its converse is false (and vice versa). If a statement's negation is false, then the statement is true (and vice versa). If a statement (or its contrapositive) and the inverse (or the converse) are both true or both false, then it is known as a logical biconditional. Traditional logic In traditional logic, contraposition is a form of immediate inference in which a proposition is inferred from another and where the former has for its subject the contradictory of the original logical proposition's predicate. In some cases, contraposition involves a change of the former's quality (i.e. affirmation or negation). For its symbolic expression in modern logic, see the rule of transposition. Contraposition also has philosophical application distinct from the other traditional inference processes of conversion and obversion where equivocation varies with different proposition types. In traditional logic, the process of contraposition is a schema composed of several steps of inference involving categorical propositions and classes. A categorical proposition contains a subject and predicate where the existential impact of the copula implies the proposition as referring to a class with at least one member, in contrast to the conditional form of hypothetical or materially implicative propositions, which are compounds of other propositions, e.g. "If P, then Q" (P and Q are both propositions), and their existential impact is dependent upon further propositions where quantification existence is instantiated (existential instantiation), not on the hypothetical or materially implicative propositions themselves. Full contraposition is the simultaneous interchange and negation of the subject and predicate, and is valid only for the type "A" and type "O" propositions of Aristotelian logic, while it is conditionally valid for "E" type propositions if a change in quantity from universal to particular is made (partial contraposition). Since the valid obverse is obtained for all the four types (A, E, I, and O types) of traditional propositions, yielding propositions with the contradictory of the original predicate, (full) contraposition is obtained by converting the obvert of the original proposition. For "E" statements, partial contraposition can be obtained by additionally making a change in quantity. Because nothing is said in the definition of contraposition with regard to the predicate of the inferred proposition, it can be either the original subject, or its contradictory, resulting in two contrapositives which are the obverts of one another in the "A", "O", and "E" type propositions. By example: from an original, 'A' type categorical proposition, All residents are voters, which presupposes that all classes have members and the existential import presumed in the form of categorical propositions, one can derive first by obversion the 'E' type proposition, No residents are non-voters. The contrapositive of the original proposition is then derived by conversion to another 'E' type proposition, No non-voters are residents. The process is completed by further obversion resulting in the 'A' type proposition that is the obverted contrapositive of the original proposition, All non-voters are non-residents. The schema of contraposition: Notice that contraposition is a valid form of immediate inference only when applied to "A" and "O" propositions. It is not valid for "I" propositions, where the obverse is an "O" proposition which has no valid converse. The contraposition of the "E" proposition is valid only with limitations (per accidens). This is because the obverse of the "E" proposition is an "A" proposition which cannot be validly converted except by limitation, that is, contraposition plus a change in the quantity of the proposition from universal to particular. Also, notice that contraposition is a method of inference which may require the use of other rules of inference. The contrapositive is the product of the method of contraposition, with different outcomes depending upon whether the contraposition is full, or partial. The successive applications of conversion and obversion within the process of contraposition may be given by a variety of names. The process of the logical equivalence of a statement and its contrapositive as defined in traditional class logic is not one of the axioms of propositional logic. In traditional logic there is more than one contrapositive inferred from each original statement. In regard to the "A" proposition this is circumvented in the symbolism of modern logic by the rule of transposition, or the law of contraposition. In its technical usage within the field of philosophic logic, the term "contraposition" may be limited by logicians (e.g. Irving Copi, Susan Stebbing) to traditional logic and categorical propositions. In this sense the use of the term "contraposition" is usually referred to by "transposition" when applied to hypothetical propositions or material implications. Form of transposition In the inferred proposition, the consequent is the contradictory of the antecedent in the original proposition, and the antecedent of the inferred proposition is the contradictory of the consequent of the original proposition. The symbol for material implication signifies the proposition as a hypothetical, or the "if–then" form, e.g. "if P, then Q". The biconditional statement of the rule of transposition (↔) refers to the relation between hypothetical (→) propositions, with each proposition including an antecedent and consequential term. As a matter of logical inference, to transpose or convert the terms of one proposition requires the conversion of the terms of the propositions on both sides of the biconditional relationship, meaning that transposing or converting to requires that the other proposition, to be transposed or converted to Otherwise, converting the terms of one proposition and not the other renders the rule invalid, violating the sufficient condition and necessary condition of the terms of the propositions, where the violation is that the changed proposition commits the fallacy of denying the antecedent or affirming the consequent by means of illicit conversion. The truth of the rule of transposition is dependent upon the relations of sufficient condition and necessary condition in logic. Sufficient condition In the proposition "If P, then Q", the occurrence of P is sufficient reason for the occurrence of Q. P, as an individual or a class, materially implicates Q, but the relation of Q to P is such that the converse proposition "If Q, then P" does not necessarily have sufficient condition. The rule of inference for sufficient condition is modus ponens, which is an argument for conditional implication: Premise (1): If P, then Q Premise (2): P Conclusion: Therefore, Q Necessary condition Since the converse of premise (1) is not valid, all that can be stated of the relationship of P and Q is that in the absence of Q, P does not occur, meaning that Q is the necessary condition for P. The rule of inference for necessary condition is modus tollens: Premise (1): If P, then Q Premise (2): not Q Conclusion: Therefore, not P Necessity and sufficiency example An example traditionally used by logicians contrasting sufficient and necessary conditions is the statement "If there is fire, then oxygen is present". An oxygenated environment is necessary for fire or combustion, but simply because there is an oxygenated environment does not necessarily mean that fire or combustion is occurring. While one can infer that fire stipulates the presence of oxygen, from the presence of oxygen the converse "If there is oxygen present, then fire is present" cannot be inferred. All that can be inferred from the original proposition is that "If oxygen is not present, then there cannot be fire". Relationship of propositions The symbol for the biconditional ("↔") signifies the relationship between the propositions is both necessary and sufficient, and is verbalized as "if and only if", or, according to the example "If P, then Q 'if and only if' if not Q, then not P". Necessary and sufficient conditions can be explained by analogy in terms of the concepts and the rules of immediate inference of traditional logic. In the categorical proposition "All S is P", the subject term S is said to be distributed, that is, all members of its class are exhausted in its expression. Conversely, the predicate term P cannot be said to be distributed, or exhausted in its expression because it is indeterminate whether every instance of a member of P as a class is also a member of S as a class. All that can be validly inferred is that "Some P are S". Thus, the type "A" proposition "All P is S" cannot be inferred by conversion from the original type "A" proposition "All S is P". All that can be inferred is the type "A" proposition "All non-P is non-S" (note that (P → Q) and (¬Q → ¬P) are both type "A" propositions). Grammatically, one cannot infer "all mortals are men" from "All men are mortal". An type "A" proposition can only be immediately inferred by conversion when both the subject and predicate are distributed, as in the inference "All bachelors are unmarried men" from "All unmarried men are bachelors". Distinguished from transposition While most authors use the terms for the same thing, some authors distinguish transposition from contraposition. In traditional logic the reasoning process of transposition as a rule of inference is applied to categorical propositions through contraposition and obversion, a series of immediate inferences where the rule of obversion is first applied to the original categorical proposition "All S is P"; yielding the obverse "No S is non-P". In the obversion of the original proposition to a type "E" proposition, both terms become distributed. The obverse is then converted, resulting in "No non-P is S", maintaining distribution of both terms. The "No non-P is S" is again obverted, resulting in the [contrapositive] "All non-P is non-S". Since nothing is said in the definition of contraposition with regard to the predicate of the inferred proposition, it is permissible that it could be the original subject or its contradictory, and the predicate term of the resulting type "A" proposition is again undistributed. This results in two contrapositives, one where the predicate term is distributed, and another where the predicate term is undistributed. Contraposition is a type of immediate inference in which from a given categorical proposition another categorical proposition is inferred which has as its subject the contradictory of the original predicate. Since nothing is said in the definition of contraposition with regard to the predicate of the inferred proposition, it is permissible that it could be the original subject or its contradictory. This is in contradistinction to the form of the propositions of transposition, which may be material implication, or a hypothetical statement. The difference is that in its application to categorical propositions the result of contraposition is two contrapositives, each being the obvert of the other, i.e. "No non-P is S" and "All non-P is non-S". The distinction between the two contrapositives is absorbed and eliminated in the principle of transposition, which presupposes the "mediate inferences" of contraposition and is also referred to as the "law of contraposition". Proof by contrapositive Because the contrapositive of a statement always has the same truth value (truth or falsity) as the statement itself, it can be a powerful tool for proving mathematical theorems (especially if the truth of the contrapositive is easier to establish than the truth of the statement itself). A proof by contrapositive is a direct proof of the contrapositive of a statement. However, indirect methods such as proof by contradiction can also be used with contraposition, as, for example, in the proof of the irrationality of the square root of 2. By the definition of a rational number, the statement can be made that "If is rational, then it can be expressed as an irreducible fraction". This statement is true because it is a restatement of a definition. The contrapositive of this statement is "If cannot be expressed as an irreducible fraction, then it is not rational". This contrapositive, like the original statement, is also true. Therefore, if it can be proven that cannot be expressed as an irreducible fraction, then it must be the case that is not a rational number. The latter can be proved by contradiction. The previous example employed the contrapositive of a definition to prove a theorem. One can also prove a theorem by proving the contrapositive of the theorem's statement. To prove that if a positive integer N is a non-square number, its square root is irrational, we can equivalently prove its contrapositive, that if a positive integer N has a square root that is rational, then N is a square number. This can be shown by setting equal to the rational expression a/b with a and b being positive integers with no common prime factor, and squaring to obtain N = a2/b2 and noting that since N is a positive integer b=1 so that N = a2, a square number. In mathematics, proof by contrapositive, or proof by contraposition, is a rule of inference used in proofs, where one infers a conditional statement from its contrapositive. In other words, the conclusion "if A, then B" is inferred by constructing a proof of the claim "if not B, then not A" instead. More often than not, this approach is preferred if the contrapositive is easier to prove than the original conditional statement itself. Logically, the validity of proof by contrapositive can be demonstrated by the use of the following truth table, where it is shown that p → q and q → p share the same truth values in all scenarios: Difference with proof by contradiction Proof by contradiction: Assume (for contradiction) that is true. Use this assumption to prove a contradiction. It follows that is false, so is true. Proof by contrapositive: To prove , prove its contrapositive statement, which is . Example Let be an integer. To prove: If is even, then is even. Although a direct proof can be given, we choose to prove this statement by contraposition. The contrapositive of the above statement is:If is not even, then is not even.This latter statement can be proven as follows: suppose that x is not even, then x is odd. The product of two odd numbers is odd, hence is odd. Thus is not even. Having proved the contrapositive, we can then infer that the original statement is true. In nonclassical logics Intuitionistic logic In intuitionistic logic, the statement cannot be proven to be equivalent to . We can prove that implies (see below) without additional assumptions, but the reverse implication, from to , requires knowing , which follows from the law of the excluded middle or an equivalent axiom. Assume (initial assumption) Assume From and , conclude Discharge assumption; conclude Turning into , conclude Discharge assumption; conclude . Subjective logic Contraposition represents an instance of the subjective Bayes' theorem in subjective logic expressed as: where denotes a pair of binomial conditional opinions given by source . The parameter denotes the base rate (aka. the prior probability) of . The pair of derivative inverted conditional opinions is denoted . The conditional opinion generalizes the logical statement , i.e. in addition to assigning TRUE or FALSE the source can assign any subjective opinion to the statement. The case where is an absolute TRUE opinion is equivalent to source saying that is TRUE, and the case where is an absolute FALSE opinion is equivalent to source saying that is FALSE. In the case when the conditional opinion is absolute TRUE the subjective Bayes' theorem operator of subjective logic produces an absolute FALSE derivative conditional opinion and thereby an absolute TRUE derivative conditional opinion which is equivalent to being TRUE. Hence, the subjective Bayes' theorem represents a generalization of both contraposition and Bayes' theorem. In probability theory Contraposition represents an instance of Bayes' theorem which in a specific form can be expressed as: In the equation above the conditional probability generalizes the logical statement , i.e. in addition to assigning TRUE or FALSE we can also assign any probability to the statement. The term denotes the base rate (aka. the prior probability) of . Assume that is equivalent to being TRUE, and that is equivalent to being FALSE. It is then easy to see that when i.e. when is TRUE. This is because so that the fraction on the right-hand side of the equation above is equal to 1, and hence which is equivalent to being TRUE. Hence, Bayes' theorem represents a generalization of contraposition.	Mathematics	Mathematical logic	null
8559544	https://en.wikipedia.org/wiki/Flexible%20shaft	Flexible shaft	A flexible shaft, often referred to as a flex shaft, is a device for transmitting rotary motion between two objects which are not fixed relative to one another. It consists of a rotating wire rope or coil which is flexible but has some torsional stiffness. It may or may not have a covering, which also bends but does not rotate. It may transmit considerable power, or only motion, with negligible power. Flexible shafts are commonly used in plumber's snakes. They are popular accessories for handheld rotary tools, and integral parts of rotary tools with a remote motor, which are called "flexible shaft tools". They are used to transmit power to some sheep shears. They are also sold to connect panel knobs to remote potentiometers or other variable electronic components. Flexible shaft tools are used frequently in the dental and jewelry industry, as well as other industrial applications.	Technology	Flexible components	null
8565007	https://en.wikipedia.org/wiki/Nonfood%20crop	Nonfood crop	A nonfood crop, also known as industrial crop, is a crop grown to produce goods for manufacturing, for example fibre for clothing, rather than food for consumption. Purpose Industrial crops is a designation given to an enterprise that attempts to raise farm sector income, and provide economic development activities for rural areas. Industrial crops also attempt to provide products that can be used as substitutes for imports from other nations. Diversity The range of crops with non-food uses is broad, but includes traditional arable crops like wheat, as well as less conventional crops like hemp and Miscanthus. Products made from non-food crops can be categorised by function:	Technology	Basics_2	null
2936554	https://en.wikipedia.org/wiki/Pont%20Alexandre%20III	Pont Alexandre III	The Pont Alexandre III () is a deck arch bridge that spans the Seine in Paris. It connects the Champs-Élysées quarter with those of the Invalides and Eiffel Tower. The bridge is widely regarded as the most ornate, extravagant bridge in the city. It has been classified as a French monument historique since 1975. History The Beaux-Arts style bridge, with its exuberant Art Nouveau lamps, cherubs, nymphs and winged horses at both ends, was built between 1896 and 1900. It is named after Tsar Alexander III of Russia, who had concluded the Franco-Russian Alliance in 1892. His son Nicholas II laid the foundation stone in October 1896. The style of the bridge reflects that of the Grand Palais, to which it leads on the right bank. The construction of the bridge is a marvel of 19th century engineering, consisting of a high single span steel arch. The design, by the architects and Gaston Cousin, was constrained by the need to keep the bridge from obscuring the view of the Champs-Élysées or the Invalides. The bridge was built by the engineers Jean Résal and . It was inaugurated in 1900 for the Exposition Universelle (universal exhibition) World's Fair, as were the nearby Grand Palais and Petit Palais. Sculptures Numerous sculptors provided the sculptures that feature prominently on the bridge. Fames Four gilt-bronze statues of Fames watch over the bridge, supported on massive masonry socles, that provide stabilizing counterweight for the arch, without interfering with monumental views. The socles are crowned by Fames restraining Pegasus. On the Right Bank: Renommée des Sciences ("Fame of the Sciences") and the Renommée des Arts ("Fame of the Arts"), both by Emmanuel Frémiet. At their bases, La France Contemporaine ("Contemporary France") by Gustave Michel and France de Charlemagne ("France of Charlemagne") by . The lions groups are by Georges Gardet. On the Left Bank: Renommée du Commerce ("Fame of Commerce") by and Renommée de l'Industrie ("Fame of Industry") by . At their bases, France de la Renaissance ("France of the Renaissance") by Jules Coutan and La France de Louis XIV ("France of Louis XIV") by Laurent Marqueste. The lions groups are by Jules Dalou. Nymphs The nymph reliefs are at the centres of the arches over the Seine, memorials to the Franco-Russian Alliance. The Nymphs of the Seine has a relief of the arms of Paris, and faces the Nymphs of the Neva with the arms of Imperial Russia. They are both executed in hammered copper over forms by Georges Récipon. In the same political spirit, the Trinity Bridge in Saint Petersburg was conceived as a memorial to the Franco-Russian Alliance. It was designed by Gustave Eiffel, and the first stone was laid in August 1897 by French president Félix Faure. Cultural associations Films and videos In the 1956 film Anastasia, the final battle takes place at the bridge. The Moody Blues' first music video footage for the song “Nights in White Satin” was shot two times with two scenes throughout the middle and the ending of the song in 1967. In the 1976 film Pardon Mon Affaire, a character drops a gun into the river from the bridge. In the 1979 film French Postcards, the final romantic scene takes place on the bridge. In the 1985 James Bond film A View to a Kill, Bond comes to a halt at the bridge in a hijacked taxi. Moments later, Bond jumps from the bridge onto a boat. In the U.S. version of "The Captain of Her Heart" music video. In the 1997 animated film Anastasia, the bridge is damaged by Rasputin in an attempt to kill Anastasia, who in real life was the granddaughter of Alexander III of Russia. Ironically, his downfall and ultimate death take place on the same bridge. In the 1998 film Ronin, the spy team meets some arms dealers under the bridge on the Right Bank. In the 2004 film A Very Long Engagement, Marion Cotillard's character kills the character played by François Levantal under the bridge. In the 2005 film Angel-A it is the Pont Alexandre III from which Angela and André jump into the Seine. In the 2006 music video for Mariah Carey's hit single "Say Somethin'" with Pharrell and Snoop Dogg. In the 2006 episode "Cold Stones" of The Sopranos, Carmela Soprano and her friend Rosalie Aprile walk in wonderment over the bridge. In the 2007 film Rush Hour 3, the taxi cab where Inspector Lee and Agent Crater rode was chased by motorcycles after escaping the casino, while a gunshot was targeted on the window of the cab. In the 2011 film Midnight in Paris, the bridge is depicted in multiple scenes, including the final one. Adele's music video for the song "Someone Like You" was shot on the bridge in 2011. In the 2016 film Me Before You, the closing shot was filmed near the northeast corner of the bridge. In the 2016 Bollywood film Befikre, the song "Nashe si Chadh Gayi" was shot on the river bank by the bridge. The 2018 film Fantastic Beasts: The Crimes of Grindelwald features a sequence with the main character, Newt Scamander, capturing an escaped magical creature known as a Zouwu on the bridge. In the 2020 Netflix Original TV series Emily in Paris, Savoir, the French marketing firm where Emily works, films a perfume advertisement here with their client, Maison Lavaux. Jung Jaehyun (NCT) was shooting a music video for his cover version of the 2017 song "I Like Me Better" by Lauv. Maison Margiela staged their Artisanal 2024 runway show in the stone chamber under the bridge. Musicals In the 2017 Broadway musical Anastasia, based on the 1997 film, the bridge is seen in the second half of the musical and in the closing scene. Anastasia was the granddaughter of Alexander III, who is mentioned in the musical. Sports In June 2017, with Paris competing against Los Angeles to host the 2024 Summer Olympics (the latter would went on to host the 2028 edition) , Paris turned some of its world-famous landmarks over to sports and installed diving boards on the Alexandre III bridge that spanned the Seine. The swimming leg of the triathlon and marathon swimming events was held here. Gallery	Technology	Bridges	null
2938620	https://en.wikipedia.org/wiki/Eucalyptol	Eucalyptol	Eucalyptol (also called cineole) is a monoterpenoid colorless liquid, and a bicyclic ether. It has a fresh camphor-like odor and a spicy, cooling taste. It is insoluble in water, but miscible with organic solvents. Eucalyptol makes up about 70–90% of eucalyptus oil. Eucalyptol forms crystalline adducts with hydrohalic acids, o-cresol, resorcinol, and phosphoric acid. Formation of these adducts is useful for purification. In 1870, F. S. Cloez identified and ascribed the name "eucalyptol" to the dominant portion of Eucalyptus globulus oil. Uses Because of its pleasant, spicy aroma and taste, eucalyptol is used in flavorings, fragrances, and cosmetics. Cineole-based eucalyptus oil is used as a flavoring at low levels (0.002%) in various products, including baked goods, confectionery, meat products, and beverages. In a 1994 report released by five top cigarette companies, eucalyptol was listed as one of the 599 additives to cigarettes. It is claimed to be added to improve the flavor. Eucalyptol is an ingredient in commercial mouthwashes, and has been used in traditional medicine as a cough suppressant. Other Eucalyptol exhibits insecticidal and insect repellent properties. In contrast, eucalyptol is one of many compounds that are attractive to males of various species of orchid bees, which gather the chemical to synthesize pheromones; it is commonly used as bait to attract and collect these bees for study. One such study with Euglossa imperialis, a nonsocial orchid bee species, has shown that the presence of cineole (also eucalyptol) elevates territorial behavior and specifically attracts the male bees. It was even observed that these males would periodically leave their territories to forage for chemicals such as cineole, thought to be important for attracting and mating with females, to synthesize pheromones. Toxicology Eucalyptol has a toxicity (LD50) of 2.48 grams per kg (rat). Ingestion in significant quantities is likely to cause headache and gastric distress, such as nausea and vomiting. Because of its low viscosity, it may directly enter the lungs if swallowed, or if subsequently vomited. Once in the lungs, it is difficult to remove and can cause delirium, convulsions, severe injury or death. Biosynthesis Eucalyptol is generated from geranyl pyrophosphate (GPP) which isomerizes to (S)-linalyl diphosphate (LPP). Ionization of the pyrophosphate, catalyzed by cineole synthase, produces eucalyptol. The process involves the intermediacy of alpha-terpinyl cation. Plants containing eucalyptol Aframomum corrorima Artemisia tridentata Cannabis Cinnamomum camphora, camphor laurel (50%) Eucalyptus globulus Eucalyptus largiflorens Eucalyptus salmonophloia Eucalyptus staigeriana Eucalyptus wandoo Hedychium coronarium, butterfly lily Helichrysum gymnocephalum Kaempferia galanga, galangal, (5.7%) S. officinalis subsp. lavandulifolia (syn. S. lavandulifolia), Spanish sage (13%) Salvia rosmarinus, rosemary Turnera diffusa, damiana Umbellularia californica, pepperwood (22.0%) Zingiber officinale, ginger	Physical sciences	Terpenes and terpenoids	Chemistry
2939202	https://en.wikipedia.org/wiki/Earth%27s%20inner%20core	Earth's inner core	Earth's inner core is the innermost geologic layer of the planet Earth. It is primarily a solid ball with a radius of about , which is about 20% of Earth's radius or 70% of the Moon's radius. There are no samples of the core accessible for direct measurement, as there are for Earth's mantle. The characteristics of the core have been deduced mostly from measurements of seismic waves and Earth's magnetic field. The inner core is believed to be composed of an iron–nickel alloy with some other elements. The temperature at its surface is estimated to be approximately , about the temperature at the surface of the Sun. The inner core is solid at high temperature because of its high pressure, in accordance with the Simon-Glatzel equation. Scientific history Earth was discovered to have a solid inner core distinct from its molten Earth's outer core in 1936, by the Danish seismologist Inge Lehmann's study of seismograms from earthquakes in New Zealand, detected by sensitive seismographs on the Earth's surface. She deduced that the seismic waves reflect off the boundary of the inner core and inferred a radius of for the inner core, not far from the currently accepted value of . In 1938, Beno Gutenberg and Charles Richter analyzed a more extensive set of data and estimated the thickness of the outer core as with a steep but continuous thick transition to the inner core, implying a radius between for the inner core. A few years later, in 1940, it was hypothesized that this inner core was made of solid iron. In 1952, Francis Birch published a detailed analysis of the available data and concluded that the inner core was probably crystalline iron. The boundary between the inner and outer cores is sometimes called the "Lehmann discontinuity", although the name usually refers to another discontinuity. The name "Bullen" or "Lehmann-Bullen discontinuity", after Keith Edward Bullen, has been proposed, but its use seems to be rare. The rigidity of the inner core was confirmed in 1971. Adam Dziewonski and James Freeman Gilbert established that measurements of normal modes of vibration of Earth caused by large earthquakes were consistent with a liquid outer core. In 2005, shear waves were detected passing through the inner core; these claims were initially controversial, but are now gaining acceptance. Data sources Seismic waves Almost all measurements that scientists have about the physical properties of the inner core are the seismic waves that pass through it. Deep earthquakes generate the most informative waves, 30 km or more below the surface of the Earth (where the mantle is relatively more homogeneous) and are recorded by seismographs as they reach the surface, all over the globe. Seismic waves include "P" (primary or pressure) compressional waves that can travel through solid or liquid materials, and "S" (secondary or shear) shear waves that can only propagate through rigid elastic solids. The two waves have different velocities and are damped at different rates as they travel through the same material. Of particular interest are the so-called "PKiKP" waves—pressure waves (P) that start near the surface, cross the mantle-core boundary, travel through the core (K), are reflected at the inner core boundary (i), cross the liquid core (K) again, cross back into the mantle, and are detected as pressure waves (P) at the surface. Also of interest are the "PKIKP" waves, that travel through the inner core (I) instead of being reflected at its surface (i). Those signals are easier to interpret when the path from source to detector is close to a straight line—namely, when the receiver is just above the source for the reflected PKiKP waves, and antipodal to it for the transmitted PKIKP waves. While S waves cannot reach or leave the inner core as such, P waves can be converted into S waves, and vice versa, as they hit the boundary between the inner and outer core at an oblique angle. The "PKJKP" waves are similar to the PKIKP waves, but are converted into S waves when they enter the inner core, travel through it as S waves (J), and are converted again into P waves when they exit the inner core. Thanks to this phenomenon, it is known that the inner core can propagate S waves, and therefore must be solid. Other sources Other sources of information about the inner core include the Earth's magnetic field. While it seems to be generated mostly by fluid and electric currents in the outer core, those currents are strongly affected by the presence of the solid inner core and by the heat that flows out of it. (Although made of iron, the core is not ferromagnetic, due to being above the Curie temperature.) the Earth's mass, its gravitational field, and its angular inertia. These are all affected by the density and dimensions of the inner layers. the natural oscillation frequencies and modes of the whole Earth oscillations, when large earthquakes make the planet "ring" like a bell. These oscillations also depend strongly on the inner layers' density, size, and shape. Physical properties Seismic wave velocity The velocity of the S waves in the core varies smoothly from about 3.7 km/s at the center to about 3.5 km/s at the surface. That is considerably less than the velocity of S waves in the lower crust (about 4.5 km/s) and less than half the velocity in the deep mantle, just above the outer core (about 7.3 km/s). The velocity of the P-waves in the core also varies smoothly through the inner core, from about 11.4 km/s at the center to about 11.1 km/s at the surface. Then the speed drops abruptly at the inner-outer core boundary to about 10.4 km/s. Size and shape On the basis of the seismic data, the inner core is estimated to be about 1221 km in radius (2442 km in diameter), which is about 19% of the radius of the Earth and 70% of the radius of the Moon. Its volume is about 7.6 billion cubic km (), which is about (0.69%) of the volume of the whole Earth. Its shape is believed to be close to an oblate ellipsoid of revolution, like the surface of the Earth, only more spherical: the flattening is estimated to be between and , meaning that the radius along the Earth's axis is estimated to be about 3 km shorter than the radius at the equator. In comparison, the flattening of the Earth as a whole is close to , and the polar radius is 21 km shorter than the equatorial one. Pressure and gravity The pressure in the Earth's inner core is slightly higher than it is at the boundary between the outer and inner cores: It ranges from about . The acceleration of gravity at the surface of the inner core can be computed to be 4.3 m/s2; which is less than half the value at the surface of the Earth (9.8 m/s2). Density and mass The density of the inner core is believed to vary smoothly from about 13.0 kg/L (= g/cm3 = t/m3) at the center to about 12.8 kg/L at the surface. As it happens with other material properties, the density drops suddenly at that surface: The liquid just above the inner core is believed to be significantly less dense, at about 12.1 kg/L. For comparison, the average density in the upper 100 km of the Earth is about 3.4 kg/L. That density implies a mass of about 1023 kg for the inner core, which is (1.7%) of the mass of the whole Earth. Temperature The temperature of the inner core can be estimated from the melting temperature of impure iron at the pressure which iron is under at the boundary of the inner core (about 330 GPa). From these considerations, in 2002, D. Alfè and others estimated its temperature as between and . However, in 2013, S. Anzellini and others obtained experimentally a substantially higher temperature for the melting point of iron, . Iron can be solid at such high temperatures only because its melting temperature increases dramatically at pressures of that magnitude (see the Clausius–Clapeyron relation). Magnetic field In 2010, Bruce Buffett determined that the average magnetic field in the liquid outer core is about 2.5 milliteslas (25 gauss), which is about 40 times the maximum strength at the surface. He started from the known fact that the Moon and Sun cause tides in the liquid outer core, just as they do on the oceans on the surface. He observed that motion of the liquid through the local magnetic field creates electric currents, that dissipate energy as heat according to Ohm's law. This dissipation, in turn, damps the tidal motions and explains previously detected anomalies in Earth's nutation. From the magnitude of the latter effect he could calculate the magnetic field. The field inside the inner core presumably has a similar strength. While indirect, this measurement does not depend significantly on any assumptions about the evolution of the Earth or the composition of the core. Viscosity Although seismic waves propagate through the core as if it were solid, the measurements cannot distinguish a solid material from an extremely viscous one. Some scientists have therefore considered whether there may be slow convection in the inner core (as is believed to exist in the mantle). That could be an explanation for the anisotropy detected in seismic studies. In 2009, B. Buffett estimated the viscosity of the inner core at 1018 Pa·s, which is a sextillion times the viscosity of water, and more than a billion times that of pitch. Composition There is still no direct evidence about the composition of the inner core. However, based on the relative prevalence of various chemical elements in the Solar System, the theory of planetary formation, and constraints imposed or implied by the chemistry of the rest of the Earth's volume, the inner core is believed to consist primarily of an iron–nickel alloy. At the estimated pressures and temperatures of the core, it is predicted that pure iron could be solid, but its density would exceed the known density of the core by approximately 3%. That result implies the presence of lighter elements in the core, such as silicon, oxygen, or sulfur, in addition to the probable presence of nickel. Recent estimates (2007) allow for up to 10% nickel and 2–3% of unidentified lighter elements. According to computations by D. Alfè and others, the liquid outer core contains 8–13% of oxygen, but as the iron crystallizes out to form the inner core the oxygen is mostly left in the liquid. Laboratory experiments and analysis of seismic wave velocities seem to indicate that the inner core consists specifically of ε-iron, a crystalline form of the metal with the hexagonal close-packed () structure. That structure can still admit the inclusion of small amounts of nickel and other elements. Structure Many scientists had initially expected that the inner core would be found to be homogeneous, because that same process should have proceeded uniformly during its entire formation. It was even suggested that Earth's inner core might be a single crystal of iron. Axis-aligned anisotropy In 1983, G. Poupinet and others observed that the travel time of PKIKP waves (P waves that travel through the inner core) was about 2 seconds less for straight north–south paths than straight paths on the equatorial plane. Even taking into account the flattening of the Earth at the poles (about 0.33% for the whole Earth, 0.25% for the inner core) and crust and upper mantle heterogeneities, this difference implied that P waves (of a broad range of wavelengths) travel through the inner core about 1% faster in the north–south direction than along directions perpendicular to that. This P wave speed anisotropy has been confirmed by later studies, including more seismic data and study of the free oscillations of the whole Earth. Some authors have claimed higher values for the difference, up to 4.8%; however, in 2017 Daniel Frost and Barbara Romanowicz confirmed that the value is between 0.5% and 1.5%. Non-axial anisotropy Some authors have claimed that P wave speed is faster in directions that are oblique or perpendicular to the N−S axis, at least in some regions of the inner core. However, these claims have been disputed by Frost and Romanowicz, who instead claim that the direction of maximum speed is as close to the Earth's rotation axis as can be determined. Causes of anisotropy Laboratory data and theoretical computations indicate that the propagation of pressure waves in the crystals of ε-iron are strongly anisotropic, too, with one "fast" axis and two equally "slow" ones. A preference for the crystals in the core to align in the north–south direction could account for the observed seismic anomaly. One phenomenon that could cause such partial alignment is slow flow ("creep") inside the inner core, from the equator towards the poles or vice versa. That flow would cause the crystals to partially reorient themselves according to the direction of the flow. In 1996, S. Yoshida and others proposed that such a flow could be caused by higher rate of freezing at the equator than at polar latitudes. An equator-to-pole flow then would set up in the inner core, tending to restore the isostatic equilibrium of its surface. Others suggested that the required flow could be caused by slow thermal convection inside the inner core. T. Yukutake claimed in 1998 that such convective motions were unlikely. However, B. Buffet in 2009 estimated the viscosity of the inner core and found that such convection could have happened, especially when the core was smaller. On the other hand, M. Bergman in 1997 proposed that the anisotropy was due to an observed tendency of iron crystals to grow faster when their crystallographic axes are aligned with the direction of the cooling heat flow. He, therefore, proposed that the heat flow out of the inner core would be biased towards the radial direction. In 1998, S. Karato proposed that changes in the magnetic field might also deform the inner core slowly over time. Multiple layers In 2002, M. Ishii and A. Dziewoński presented evidence that the solid inner core contained an "innermost inner core" (IMIC) with somewhat different properties than the shell around it. The nature of the differences and radius of the IMIC are still unresolved as of 2019, with proposals for the latter ranging from 300 km to 750 km. A. Wang and X. Song proposed, in 2018, a three-layer model, with an "inner inner core" (IIC) with about 500 km radius, an "outer inner core" (OIC) layer about 600 km thick, and an isotropic shell 100 km thick. In this model, the "faster P wave" direction would be parallel to the Earth's axis in the OIC, but perpendicular to that axis in the IIC. However, the conclusion has been disputed by claims that there need not be sharp discontinuities in the inner core, only a gradual change of properties with depth. In 2023, a study reported new evidence "for an anisotropically-distinctive innermost inner core" – a ~650-km thick innermost ball – "and its transition to a weakly anisotropic outer shell, which could be a fossilized record of a significant global event from the past." They suggest that atoms in the IIC atoms are [packed] slightly differently than its outer layer, causing seismic waves to pass through the IIC at different speeds than through the surrounding core (P-wave speeds ~4% slower at ~50° from the Earth’s rotation axis). Lateral variation In 1997, S. Tanaka and H. Hamaguchi claimed, on the basis of seismic data, that the anisotropy of the inner core material, while oriented N−S, was more pronounced in "eastern" hemisphere of the inner core (at about 110 °E longitude, roughly under Borneo) than in the "western" hemisphere (at about 70 °W, roughly under Colombia). Alboussère and others proposed that this asymmetry could be due to melting in the Eastern hemisphere and re-crystallization in the Western one. C. Finlay conjectured that this process could explain the asymmetry in the Earth's magnetic field. However, in 2017 Frost and Romanowicz disputed those earlier inferences, claiming that the data shows only a weak anisotropy, with the speed in the N−S direction being only 0.5% to 1.5% faster than in equatorial directions, and no clear signs of E−W variation. Other structure Other researchers claim that the properties of the inner core's surface vary from place to place across distances as small as 1 km. This variation is surprising since lateral temperature variations along the inner-core boundary are known to be extremely small (this conclusion is confidently constrained by magnetic field observations). Growth The Earth's inner core is thought to be slowly growing as the liquid outer core at the boundary with the inner core cools and solidifies due to the gradual cooling of the Earth's interior (about 100 degrees Celsius per billion years). According to calculations by Alfé and others, as the iron crystallizes onto the inner core, the liquid just above it becomes enriched in oxygen, and therefore less dense than the rest of the outer core. This process creates convection currents in the outer core, which are thought to be the prime driver for the currents that create the Earth's magnetic field. The existence of the inner core also affects the dynamic motions of liquid in the outer core, and thus may help fix the magnetic field. Dynamics Because the inner core is not rigidly connected to the Earth's solid mantle, the possibility that it rotates slightly more quickly or slowly than the rest of Earth has long been entertained. In the 1990s, seismologists made various claims about detecting this kind of super-rotation by observing changes in the characteristics of seismic waves passing through the inner core over several decades, using the aforementioned property that it transmits waves more quickly in some directions. In 1996, X. Song and P. Richards estimated this "super-rotation" of the inner core relative to the mantle as about one degree per year. In 2005, they and J. Zhang compared recordings of "seismic doublets" (recordings by the same station of earthquakes occurring in the same location on the opposite side of the Earth, years apart), and revised that estimate to 0.3 to 0.5 degree per year. In 2023, it was reported that the core stopped spinning faster than the planet's surface around 2009 and likely is now rotating slower than it. This is not thought to have major effects and one cycle of the oscillation is thought to be about seven decades, coinciding with several other geophysical periodicities, "especially the length of day and magnetic field". In 1999, M. Greff-Lefftz and H. Legros noted that the gravitational fields of the Sun and Moon that are responsible for ocean tides also apply torques to the Earth, affecting its axis of rotation and a slowing down of its rotation rate. Those torques are felt mainly by the crust and mantle, so that their rotation axis and speed may differ from overall rotation of the fluid in the outer core and the rotation of the inner core. The dynamics is complicated because of the currents and magnetic fields in the inner core. They find that the axis of the inner core wobbles (nutates) slightly with a period of about 1 day. With some assumptions on the evolution of the Earth, they conclude that the fluid motions in the outer core would have entered resonance with the tidal forces at several times in the past (3.0, 1.8, and 0.3 billion years ago). During those epochs, which lasted 200–300 million years each, the extra heat generated by stronger fluid motions might have stopped the growth of the inner core. Age Theories about the age of the core are part of theories of the history of Earth. It is widely believed that the Earth's solid inner core formed out of an initially completely liquid core as the Earth cooled. However, the time when this process started is unknown. Two main approaches have been used to infer the age of the inner core: thermodynamic modeling of the cooling of the Earth, and analysis of paleomagnetic evidence. The estimates yielded by these methods vary from 0.5 to 2 billion years old. Thermodynamic evidence One of the ways to estimate the age of the inner core is by modeling the cooling of the Earth, constrained by a minimum value for the heat flux at the core–mantle boundary (CMB). That estimate is based on the prevailing theory that the Earth's magnetic field is primarily triggered by convection currents in the liquid part of the core, and the fact that a minimum heat flux is required to sustain those currents. The heat flux at the CMB at present time can be reliably estimated because it is related to the measured heat flux at Earth's surface and to the measured rate of mantle convection. In 2001, S. Labrosse and others, assuming that there were no radioactive elements in the core, gave an estimate of 1±0.5 billion years for the age of the inner core — considerably less than the estimated age of the Earth and of its liquid core (about 4.5 billion years) In 2003, the same group concluded that, if the core contained a reasonable amount of radioactive elements, the inner core's age could be a few hundred million years older. In 2012, theoretical computations by M. Pozzo and others indicated that the electrical conductivity of iron and other hypothetical core materials, at the high pressures and temperatures expected there, were two or three times higher than assumed in previous research. These predictions were confirmed in 2013 by measurements by Gomi and others. The higher values for electrical conductivity led to increased estimates of the thermal conductivity, to 90 W/m·K; which, in turn, lowered estimates of its age to less than 700 million years old. However, in 2016 Konôpková and others directly measured the thermal conductivity of solid iron at inner core conditions, and obtained a much lower value, 18–44 W/m·K. With those values, they obtained an upper bound of 4.2 billion years for the age of the inner core, compatible with the paleomagnetic evidence. In 2014, Driscoll and Bercovici published a thermal history of the Earth that avoided the so-called mantle thermal catastrophe and new core paradox by invoking 3 TW of radiogenic heating by the decay of in the core. Such high abundances of K in the core are not supported by experimental partitioning studies, so such a thermal history remains highly debatable. Paleomagnetic evidence Another way to estimate the age of the Earth is to analyze changes in the magnetic field of Earth during its history, as trapped in rocks that formed at various times (the "paleomagnetic record"). The presence or absence of the solid inner core could result in different dynamic processes in the core that could lead to noticeable changes in the magnetic field. In 2011, Smirnov and others published an analysis of the paleomagnetism in a large sample of rocks that formed in the Neoarchean (2.8–2.5 billion years ago) and the Proterozoic (2.5–0.541 billion). They found that the geomagnetic field was closer to that of a magnetic dipole during the Neoarchean than after it. They interpreted that change as evidence that the dynamo effect was more deeply seated in the core during that epoch, whereas in the later time currents closer to the core-mantle boundary grew in importance. They further speculate that the change may have been due to growth of the solid inner core between 3.5–2.0 billion years ago. In 2015, Biggin and others published the analysis of an extensive and carefully selected set of Precambrian samples and observed a prominent increase in the Earth's magnetic field strength and variance around 1.0–1.5 billion years ago. This change had not been noticed before due to the lack of sufficient robust measurements. They speculated that the change could be due to the birth of Earth's solid inner core. From their age estimate they derived a rather modest value for the thermal conductivity of the outer core, that allowed for simpler models of the Earth's thermal evolution. In 2016, P. Driscoll published a numerical evolving dynamo model that made a detailed prediction of the paleomagnetic field evolution over 0.0–2.0 Ga. The evolving dynamo model was driven by time-variable boundary conditions produced by the thermal history solution in Driscoll and Bercovici (2014). The evolving dynamo model predicted a strong-field dynamo prior to 1.7 Ga that is multipolar, a strong-field dynamo from 1.0–1.7 Ga that is predominantly dipolar, a weak-field dynamo from 0.6–1.0 Ga that is a non-axial dipole, and a strong-field dynamo after inner core nucleation from 0.0–0.6 Ga that is predominantly dipolar. An analysis of rock samples from the Ediacaran epoch (formed about 565 million years ago), published by Bono and others in 2019, revealed unusually low intensity and two distinct directions for the geomagnetic field during that time that provides support for the predictions by Driscoll (2016). Considering other evidence of high frequency of magnetic field reversals around that time, they speculate that those anomalies could be due to the onset of formation of the inner core, which would then be 0.5 billion years old. A News and Views by P. Driscoll summarizes the state of the field following the Bono results. New paleomagnetic data from the Cambrian appear to support this hypothesis.	Physical sciences	Geophysics	null
2940782	https://en.wikipedia.org/wiki/Notonectidae	Notonectidae	Notonectidae is a cosmopolitan family of aquatic insects in the order Hemiptera, commonly called backswimmers because they swim "upside down" (inverted). They are all predators and typically range from in length. They are similar in appearance to Corixidae (water boatmen), but can be separated by differences in their dorsal-ventral coloration, front legs, and predatory behavior. Their dorsum is convex, lightly colored without cross striations. Their front tarsi are not scoop-shaped and their hind legs are fringed for swimming. There are about 350 species in two subfamilies: Notonectinae with seven genera, and Anisopinae with four genera. Members in the former subfamily are often larger than those in the latter. Backswimmers swim on their backs, vigorously paddling with their long, hair-fringed hind legs and attack prey as large as tadpoles and small fish. They can inflict a painful "bite" on a human being, actually a stab with their sharp tubular mouthparts (proboscis). They inhabit still freshwater, e.g. lakes, ponds, marshes, and are sometimes found in garden ponds and even swimming pools. Although primarily aquatic, they can fly well and so can disperse easily to new habitats. The best-known genus of backswimmers is Notonecta – streamlined, deep-bodied bugs up to long, green, brown, or yellowish in colour. The common backswimmer, N. glauca, is widespread in Europe, including the United Kingdom where it is known as the greater water boatman. Another of the same region, N. maculata, is distinguished by its mottled brick-coloured forewings. In contrast to other aquatic insects that cling to submerged objects, the two genera Anisops and Buenoa uses a unique system to stay submerged: using the extra oxygen supply from haemoglobin in their abdomen, instead of using oxygen dissolved in the water. The size of these air bubbles, which provide buoyancy, changes as the nitrogen dissolves into the blood and the oxygen is used in respiration. This allows for regulation of the size of the air bubbles and their concentration of oxygen.	Biology and health sciences	Hemiptera (true bugs)	Animals
2942638	https://en.wikipedia.org/wiki/Gloss%20%28optics%29	Gloss (optics)	Gloss is an optical property which indicates how well a surface reflects light in a specular (mirror-like) direction. It is one of the important parameters that are used to describe the visual appearance of an object. Other categories of visual appearance related to the perception of regular or diffuse reflection and transmission of light have been organized under the concept of cesia in an order system with three variables, including gloss among the involved aspects. The factors that affect gloss are the refractive index of the material, the angle of incident light and the surface topography. Apparent gloss depends on the amount of specular reflection – light reflected from the surface in an equal amount and the symmetrical angle to the one of incoming light – in comparison with diffuse reflection – the amount of light scattered into other directions. Theory When light illuminates an object, it interacts with it in a number of ways: Absorbed within it (largely responsible for colour) Transmitted through it (dependent on the surface transparency and opacity) Scattered from or within it (diffuse reflection, haze and transmission) Specularly reflected from it (gloss) Variations in surface texture directly influence the level of specular reflection. Objects with a smooth surface, i.e. highly polished or containing coatings with finely dispersed pigments, appear shiny to the eye due to a large amount of light being reflected in a specular direction whilst rough surfaces reflect no specular light as the light is scattered in other directions and therefore appears dull. The image forming qualities of these surfaces are much lower making any reflections appear blurred and distorted. Substrate material type also influences the gloss of a surface. Non-metallic materials, i.e. plastics etc. produce a higher level of reflected light when illuminated at a greater illumination angle due to light being absorbed into the material or being diffusely scattered depending on the colour of the material. Metals do not suffer from this effect producing higher amounts of reflection at any angle. The Fresnel formula gives the specular reflectance, , for an unpolarized light of intensity , at angle of incidence , giving the intensity of specularly reflected beam of intensity , while the refractive index of the surface specimen is . The Fresnel equation is given as follows : Surface roughness Surface roughness influences the specular reflectance levels; in the visible frequencies, the surface finish in the micrometre range is most relevant. The diagram on the right depicts the reflection at an angle on a rough surface with a characteristic roughness height variation . The path difference between rays reflected from the top and bottom of the surface bumps is: When the wavelength of the light is , the phase difference will be: If is small, the two beams (see Figure 1) are nearly in phase, resulting in constructive interference; therefore, the specimen surface can be considered smooth. But when , then beams are not in phase and through destructive interference, cancellation of each other will occur. Low intensity of specularly reflected light means the surface is rough and it scatters the light in other directions. If the middle phase value is taken as criterion for smooth surface, , then substitution into the equation above will produce: This smooth surface condition is known as the Rayleigh roughness criterion. History The earliest studies of gloss perception are attributed to Leonard R. Ingersoll who in 1914 examined the effect of gloss on paper. By quantitatively measuring gloss using instrumentation Ingersoll based his research around the theory that light is polarised in specular reflection whereas diffusely reflected light is non-polarized. The Ingersoll "glarimeter" had a specular geometry with incident and viewing angles at 57.5°. Using this configuration gloss was measured using a contrast method which subtracted the specular component from the total reflectance using a polarizing filter. In the 1930s work by A. H. Pfund, suggested that although specular shininess is the basic (objective) evidence of gloss, actual surface glossy appearance (subjective) relates to the contrast between specular shininess and the diffuse light of the surrounding surface area (now called "contrast gloss" or "luster"). If black and white surfaces of the same shininess are visually compared, the black surface will always appear glossier because of the greater contrast between the specular highlight and the black surroundings as compared to that with white surface and surroundings. Pfund was also the first to suggest that more than one method was needed to analyze gloss correctly. In 1937 R. S. Hunter, as part of his research paper on gloss, described six different visual criteria attributed to apparent gloss. The following diagrams show the relationships between an incident beam of light, I, a specularly reflected beam, S, a diffusely reflected beam, D and a near-specularly reflected beam, B. Specular gloss – the perceived brightness and the brilliance of highlights Defined as the ratio of the light reflected from a surface at an equal but opposite angle to that incident on the surface. Sheen – the perceived shininess at low grazing angles Defined as the gloss at grazing angles of incidence and viewing Contrast gloss – the perceived brightness of specularly and diffusely reflecting areas Defined as the ratio of the specularly reflected light to that diffusely reflected normal to the surface; Absence of bloom – the perceived cloudiness in reflections near the specular direction Defined as a measure of the absence of haze or a milky appearance adjacent to the specularly reflected light: haze is the inverse of absence-of-bloom Distinctness of image gloss – identified by the distinctness of images reflected in surfaces Defined as the sharpness of the specularly reflected light Surface texture gloss – identified by the lack of surface texture and surface blemishes Defined as the uniformity of the surface in terms of visible texture and defects (orange peel, scratches, inclusions etc.) A surface can therefore appear very shiny if it has a well-defined specular reflectance at the specular angle. The perception of an image reflected in the surface can be degraded by appearing unsharp, or by appearing to be of low contrast. The former is characterised by the measurement of the distinctness-of-image and the latter by the haze or contrast gloss. In his paper Hunter also noted the importance of three main factors in the measurement of gloss: The amount of light reflected in the specular direction The amount and way in which the light is spread around the specular direction The change in specular reflection as the specular angle changes For his research he used a glossmeter with a specular angle of 45° as did most of the first photoelectric methods of that type, later studies however by Hunter and D. B. Judd in 1939, on a larger number of painted samples, concluded that the 60 degree geometry was the best angle to use so as to provide the closest correlation to a visual observation. Standard gloss measurement Standardisation in gloss measurement was led by Hunter and ASTM (American Society for Testing and Materials) who produced ASTM D523 Standard test method for specular gloss in 1939. This incorporated a method for measuring gloss at a specular angle of 60°. Later editions of the Standard (1951) included methods for measuring at 20° for evaluating high gloss finishes, developed at the DuPont Company (Horning and Morse, 1947) and 85° (matte, or low, gloss). ASTM has a number of other gloss-related standards designed for application in specific industries including the old 45° method which is used primarily now used for glazed ceramics, polyethylene and other plastic films. In 1937, the paper industry adopted a 75° specular-gloss method because the angle gave the best separation of coated book papers. This method was adopted in 1951 by the Technical Association of Pulp and Paper Industries as TAPPI Method T480. In the paint industry, measurements of the specular gloss are made according to International Standard ISO 2813 (BS 3900, Part 5, UK; DIN 67530, Germany; NFT 30-064, France; AS 1580, Australia; JIS Z8741, Japan, are also equivalent). This standard is essentially the same as ASTM D523 although differently drafted. Studies of polished metal surfaces and anodised aluminium automotive trim in the 1960s by Tingle, Potter and George led to the standardisation of gloss measurement of high gloss surfaces by goniophotometry under the designation ASTM E430. In this standard it also defined methods for the measurement of distinctness of image gloss and reflection haze.	Physical sciences	Basics_7	null
1484327	https://en.wikipedia.org/wiki/Guangzhou%20Metro	Guangzhou Metro	The Guangzhou Metro () is the rapid transit system of the city of Guangzhou in the Guangdong Province of China. It is operated by the state-owned Guangzhou Metro Corporation and was the fourth metro system to be built in mainland China, after those of Beijing, Tianjin, and Shanghai. The earliest efforts to build an underground rapid transit system in Guangzhou date back to 1960. In the two decades that followed, the project was brought into the agenda five times but ended up abandoned each time due to financial and technical difficulties. Preparation of what would lead to today's Guangzhou Metro did not start until the 1980s, and it was not until 1993 that construction of the first line, Line 1, officially began. Line 1 opened four years later in 1997 with five stations in operation. , Guangzhou Metro has 17 lines in operation, namely: Line 1, Line 2, Line 3, Line 4, Line 5, Line 6, Line 7, Line 8, Line 9, Line 11, Line 13, Line 14, Line 18, Line 21, Line 22, Guangfo Line, and Zhujiang New Town APM reaching both the urban core and surrounding suburbs. Guangfo Line connects Guangzhou and Foshan and is the first metro line between two cities in the country. Daily service hours start at 6:00 am and end at midnight and daily ridership averages over 7 million. Having delivered 3.029 billion rides in 2018, Guangzhou Metro is the third busiest metro system in the world and the 3rd largest in terms of length, after the metro systems of Beijing and Shanghai. Guangzhou Metro operates 320 stations and of lines. Extensive development of the metro network has been planned for the next decade, with construction started on Line 10, Line 12, and Line 24, and extensions of Line 8, Line 13, Line 14, Line 18, as well as the extension of Line 22 to Baiyun Airport. Some of the system's lines were designed to operate much faster than traditional metro lines, with stations far apart and faster trainsets regularly running at . Lines 18 and 22 are the fastest metro lines in China, a title previously held by Line 11 of the Shenzhen Metro. History Forays of the 1960s and 1970s Chen Yu (), Governor of Guangdong in 1957–1967, was the first to have proposed an underground metro system for Guangzhou. In the summer of 1960, he ordered a secret geological survey of groundwater levels of Guangzhou. Six holes with an accumulated depth of were drilled in the karst and alluvial plains in the city. The geological conditions of Guangzhou, despite their complexity, did not preclude the possibility of an underground metro system. Analysis of the survey data resulted in a confidential report titled Geological Survey for Guangzhou Underground Railway Project dated July 1961, the earliest one of such reports. In 1965, Chen Yu along with Tao Zhu (), who had been the Governor of Guangdong and First Secretary of Guangdong Committee of the Chinese Communist Party, proposed in the wake of the Gulf of Tonkin incident that a tunnel is built in Guangzhou for wartime evacuations and post-war metro development. Approved by the central government, the project started in the spring of 1965. Due to its confidentiality in the context of intensification of the Vietnam War, the project adopted the obscure name of "Project Nine" (), where "Nine" was the number of strokes in "", the Chinese word for "underground". As envisaged by Chen Yu, the metro system of Guangzhou would consist of two lines: a north–south line that would connect Nanfang Building to Sanyuanli via Renmin Lu and Jiefang Beilu, and an east–west line that would run from Xichang to Dongshan along today's Dongfeng Lu. The two lines roughly parallelled Line 2 and Line 1 of the modern days, respectively. The east–west line was never built, while Project Nine was dedicated to the north–south line. Over ten teams of miners were recruited for a project filled with hazards and perils. Constrained by extreme scarcity of time, monetary and material resources, the ambition to build a tunnel for the metro operation was scaled back— the capability to run trolleybuses was deemed acceptable. For ¥13 million, an long tunnel was completed in 1966. The tunnel was planned to be used as an air-raid shelter and eventual metro line; however, with a cross-section merely 3 m wide and 2.85 m tall, and exposed rocks and wooden trestles scattered everywhere, it was unusable for public transit. In the two decades that followed, four attempts were made to revive and expand Project Nine, first in 1970, next in 1971, then in 1974, and last in 1979. Due to lack of funds and complex geotechnical conditions, none of these efforts materialized. Construction of Line 1 The metro project of Guangzhou was launched for the sixth time in 1984 as the Preparation Office of Guangzhou Metro, established back in 1979 as part of the last attempt to resurrect Project Nine, was moved out of the civil air-defense system and became a subordinate body of the Construction Commission of Guangzhou, bringing Guangzhou Metro into the scope of urban infrastructure development. Before the 1980s, war preparedness was the dominant tenet of underground infrastructure projects in mainland China. The construction of Guangzhou Metro marked the first deviation from the old doctrine as traffic itself became the prime consideration of the project. The design of the initial metro network was a collaborative effort between China and France (SYSTRA). Four tentative designs were published on 14 March 1988 edition of Guangzhou Daily. From the four designs, one was selected based on expert and mass feedback. The selected design, featuring two intersecting lines, was the baseline typology for today's Line 1 and Line 2. Construction of Line 1 officially commenced on 28 December 1993, although work on a trial section at Huangsha had begun in October 1992, five months before the feasibility study of the line was ratified by the State Planning Commission in March 1993. Various technologies novel to China's construction industry at the time were adopted in different sections of the project, notably including immersed tubes (Pearl River Tunnel) and tunnel boring machines (Huangsha–Martyrs' Park section). As the most massive urban infrastructure project in the history of Guangzhou, Line 1 required funding of ¥12.75 billion, all of which was raised by the local government. Use of cut-and-cover tunnels aggressively backed by then-mayor Li Ziliu necessitated the relocation of approximately 100,000 residents in 20,000 households and demolition of buildings totalling in the area and earned Li the nickname "Li the Demolisher" (). Three and a half years after construction started, the section from Xilang to Huangsha opened for trial operation on 28 June 1997. The remaining , from Huangsha to Guangzhou East railway station, was completed eighteen months later on 28 December 1998. The entire line opened for sightseeing tours between 16 February and 2 March 1999, delivering 1.39 million rides 15 days before closing for final testing. Operation of Line 1 officially began on 28 June 1999, 34 years after the start of Project Nine in 1965. Accelerated expansion in the 2000s The success of Line 1 as a turnkey project acquired from Siemens with 100% imported electromechanical equipment prompted a wave of similar proposals from twelve other cities in mainland China toward the end of the 1990s. The fever for import-centric rapid transit caused the State Planning Committee to temporarily halt approval of rapid transit projects nationwide and regulate the localization rates of rolling stock suppliers. Amid tightened regulation, only Line 2 of Guangzhou Metro received the immediate green light to proceed in June 1998 on the condition that at least 60% of its electromechanical equipment must be sourced domestically. Construction of Line 2 started in July 1998. Rolling stock manufacturer Bombardier airlifted the first two train cars in an An-124 from Berlin to Guangzhou in November 2002 after schedule delays. The first section, from to opened on 29 December 2002; the remaining section from Xiaogang to opened on 28 June 2003. At ¥2.13 billion, the equipment cost of Line 2 was 53% lower than that of Line 1. This demonstrated the feasibility of cost reduction through procurement of domestic equipment, revealing a path to project approval to other Chinese cities and reigniting their aspirations to own a rapid transit system. The renewed craze for rapid transit across the country soon encountered a new round of tightened control on project approval around 2003. But Guangzhou was exempted along with Beijing, Shanghai and Shenzhen. By the time Line 2 was completed, construction of Line 3, Line 4, and Guangfo Line had been underway, among which only Guangfo Line later fell to stringent regulation of approvals. Lines in operation Line 1 Line 1 runs from Xilang to Guangzhou East railway station, with a total length of . Except for Kengkou and Xilang, all stations in Line 1 are underground. Its first section, from Xilang to Huangsha, opened on 28 June 1997, making Guangzhou the fourth city in mainland China to have a metro system. The full line started operation two years later on 28 June 1999. Line 1's color is yellow. Line 2 Line 2 is a north–south line that runs from Jiahewanggang to Guangzhou South railway station. Until 21 September 2010, it ran from to Wanshengwei. Its first section, between Sanyuanli and , opened on 29 December 2002. It was extended from Xiaogang to on 28 June 2003 and further to Wanshengwei a year later. The section between Xiaogang and Wanshengwei was split off to form part of Line 8 during 22–24 September 2010, when the operation was paused. The latest extension, from to Guangzhou South railway station and from to , opened on 25 September 2010 as the whole line resumed operation. The length of the current line is . All stations in Line 2 are underground. Line 2's color is deep blue. Line 3 Line 3 is a Y-shaped line connecting Airport North and Tianhe Coach Terminal to Haibang. All stations in the line are underground. When the line opened on 26 December 2005, trains operated between Guangzhou East railway station and Kecun. Following completion of the Tianhe Coach Terminal–Tiyu Xilu and Kecun–Panyu Square sections, the line was rerouted on 30 December 2006 to offer transfer-free connections between Panyu Square and Tianhe Coach Terminal via Tiyu Xilu. The Guangzhou East railway station–Tiyu Xilu section became a shuttle until it was extended northwards to Airport South on 30 October 2010. Southwards, it was extended from Panyu Square to Haibang on 1 November 2024. In official distinctions, the main route consists of the entire Airport North–Haibang section, while the Tianhe Coach Terminal–Tiyu Xilu section is a spur line. The spur line will be split off in the long term to form part of Line 10. Line 3 had been notorious for its crowding since it opened, for it ran three-car trains. That was partly relieved when all three-car trains started operating as six-car ones, connected in sets of two, on 28 April 2010. Sectional services between Tonghe to Dashi are added from 7:30 to 8:30 every workday, partly solving the capacity issues. Despite these changes, as of 2018, the line is still severely overcrowded. Line 3's color is orange. Line 4 Line 4 is a north–south line running parallel to Line 2 along the east of the city. It is long with 24 stations. The section of the line from Huangcun to Xinzao, Feishajiao to Nansha Passenger Port are built underground, while that from Xinzao to Jinzhou is built at the elevated track. It was the first metro line in mainland China to use linear motor trains. Its first section, from Wanshengwei to Xinzao, opened on 26 December 2005. Southwards, it was extended from Xinzao to Huangge on 30 December 2006 and further to Jinzhou on 28 June 2007. Northwards, it was extended to Chebeinan on 28 December 2009. Southwards, it extended from Chebeinan to Huangcun, opened on 25 September 2010. Its latest extension, from Huangcun to Nansha Passenger Port, opened on 27 December 2017. Line 4's color is green. Line 5 The long Line 5 starts at Jiaokou and runs to Huangpu New Port. It entered operation on 28 December 2009 between Jiaokou and Wenchong, and on 28 December 2023 between Wenchong and Huangpu New Port. All stations in the line except Jiaokou and Tanwei are underground. Until Line 8 was split off from Line 2, it was the only line that interchanged with all other lines. Similar to Line 4, Line 5 also uses linear motor trains. Line 5's color is red. Line 6 The first stage of Line 6, a long phase one runs from Xunfenggang to Changban with 22 stations. It began service on 28 December 2013 and contains three elevated stations along the route. Construction of a 10-station, long extension to Xiangxue from Changban is entered revenue service in 2016. The line runs four-car trains, but stations of the east extension starting with South China Botanical Garden will be constructed with a provision to accommodate six-car trains in preparation for a route split in the future. Line 6's color is maroon. Line 7 The first phase of Line 7 began service on 28 December 2016 and ran from Guangzhou South railway station to Higher Education Mega Center South in Panyu District throughout . The phase 1 west extension opened on 1 May 2022 from Guangzhou South railway station to Meidi Dadao station. Six-car trains are used. All stations are underground. Phase 2 opened on 28 December 2023, and extends the line by and 11 stations to reach north of the Pearl River and go deep to Huangpu district, providing interchanges with Line 13 at , Line 5 at , Line 6 at , Line 21 at and the planned east extension of Line 8 at . Line 7's color is light green. Line 8 The first section of Line 8, from Xiaogang to Wanshengwei, opened in 2002 and ran as part of Line 2 until the extension to the line was completed in September 2010. Line 8 ran from Fenghuang Xincun to Wanshengwei. The section from Changgang to Wanshengwei opened on 25 September 2010 when the split-off from Line 2 was complete. The section west of Changgang did not open until 3 November 2010 due to disputes over the environmental impact of the cooling facilities at Shayuan. The remaining section from Fenghuang Xincun to Cultural Park and Cultural Park to Jiaoxin are opened on 28 December 2019 and 26 November 2020 separately. Line 8's color is teal. Line 9 The long underground route is operated by six-car trains, which runs from Fei'eling to Gaozeng, serving 10 stations. The line, other than Qingtang station, went operational on 28 December 2017. Line 9 mainly serves as a link for the passengers of Huadu District and Guangzhou North railway station to the rest of the system, having only one transfer station with Line 3 at Gaozeng. After the Tianhe Coach Terminal–Tiyu Xilu spur line of Line 3 is split off to form part of Line 10, the line is expected to be connected into Line 3 using the reserved switches at Gaozeng to become a new spur line. Line 9's color is pale green. Line 11 Line 11 is a loop-shaped line—the first in such shape—connecting and , via Guangzhou railway station, Guangzhou East railway station, , , and . The line was opened on 28 December 2024 at 14:00 local time, with trains stopping in all but the Guangzhou East and s. Line 11's color is gold. Line 13 Opened on 28 December 2017, Line 13 is the first metro line in Guangzhou built to run eight-car trains. The currently operating first phase runs from Yuzhu to Xinsha, serving passengers of Huangpu and Xintang, Zengcheng. The eleven-station line currently has only one transfer station with Line 5 at Yuzhu. The second phase of Line 13 runs west of the current phase, which cuts through popular areas of Huangpu, Tianhe, and Liwan Districts, and is currently under construction. Line 13's color is olive. Line 14 Two sections of Line 14 are currently in service. The Knowledge City Branch Line, a ten-station long route located mainly within Huangpu, opened on 28 December 2017. The branch line operates primarily within Huangpu between Xinhe and Zhenlong, serving the Sino-Singapore Guangzhou Knowledge City. The mainline segment to Conghua opened a year later on 28 December 2018 and runs from Jiahewanggang in Baiyun District to Dongfeng in Conghua. A southward extension to Guangzhou railway station is currently under construction. Line 14 was the first Guangzhou Metro line to run express services. Line 14's color is brown. Line 18 The section from to of Line 18 opened on 28 September 2021. The section is 58.3 km in length. It will be extended 3 km to . A further 39.6 km extension to is also planned. Line 18's color is blue. Line 21 The long Line 21 runs between Tianhe Park in Tianhe and Zengcheng Square in Zengcheng with six-car trains. It has of underground tracks, of elevated tracks, and of tracks in mountain tunnels. The section from Yuancun to Tianhe Park is intended as part of Line 11 and constructed to accommodate the eight-car trains of the latter. When the construction of Line 11 is completed, this section will be operated as part of Line 11, making Tianhe Park the west end of Line 21. Express service was also provided after the inauguration of the western section (Yuancun – Zhenlongxi). Line 21's color is dark navy. Line 22 The section from to of Line 22 opened on 31 March 2022. The section is 18.2 km in length. It will be extended 73.2 km to . Line 22's color is orange. Guangfo Line The Guangzhou–Foshan Section of Pearl River Delta Region Intercity Rapid Transit () is an intercity metro line that connects Guangzhou and Foshan. It is commonly known as Guangfo Metro and Guangfo Line of Guangzhou Metro. The section within Foshan also doubles as Line 1 of FMetro (Foshan Metro). The line is operated by Guangdong Guangfo Rail Transit Co., Ltd., a subsidiary co-owned by Guangzhou Metro (51%) and Foshan Metro (49%). Its first section, from Xilang to Kuiqi Lu in Foshan, started operation on 3 November 2010 with of tracks and 14 stations. Eleven of the stations are located in Foshan, while the other three are in Guangzhou. Relocation disputes at Lijiao were not resolved until October 2013 and have delayed completion of the extension from Xilang to Lijiao till December 2015. When the line is completed, it will have of tracks and 21 stations, of which of tracks and 10 stations will be located in Guangzhou. The line runs four-car trains. All its stations are underground. Zhujiang New Town APM Line The Automated People Mover System of Zhujiang New Town Core District Municipal Traffic Project () is an underground automated people mover that serves the central business district of Zhujiang New Town. It is commonly known as Zhujiang New Town Automated People Mover System or the APM for short. At a length of , it connects Linhexi and Canton Tower with nine stations on the line. The operation started on 8 November 2010 with Canton Tower Station named Chigang Pagoda Station until December 2013. The stations of Haixinsha and Chigang Pagoda remained closed during the 2010 Asian Games. Chigang Pagoda Station opened on 28 November 2010, one day after the Asian Games ended; Haixinsha Station remained unopened until 24 February 2011. There is no direct platform-to-platform connection between the APM and Line 3 albeit they share the stations of Linhexi and Canton Tower. Transfer passengers need to exit and reenter with a new ticket. The APM runs two-car rubber-wheeled driverless trains. Network expansion Short-term planning Long-term planning The Guangzhou Urban Rail Transit Network Planning Scheme (2018–2035) (), which was approved by the Guangzhou Municipal Government in November 2020, shows that a total of 53 metro lines and 2,029 km are planned in Guangzhou. This round of line network planning is divided into three levels: high-speed metro, rapid metro, and regular-speed metro. Among them, there are 5 high-speed metro lines with 452 km in Guangzhou, 11 rapid metro lines with 607 km in Guangzhou, and 37 regular-speed metro lines with 970 km. High-speed metro lines: : Knowledge City – Luogang – Zini (→ Foshan) : (Zhongshan / Zhuhai →) Shiliuchong – Huachengjie (→ Qingyuan) : Airport North – Nansha Passenger Port (→ Dongguan) : (Foshan →) Fangcun – Xintang (→ Dongguan) spur line: Xintang – Guangzhou Huali College (→ Huizhou) : Guangzhou East railway station – Liangkou (→ Xinfeng) Rapid metro lines: : Airport North – Haiou Island parallel express line: Pazhou – Jiaomen : Shuixibei – Meidi Dadao : Chaoyang – Xinsha : Guangzhou railway station – Dongfeng : Xintang – Lichengbei : Tianhe Park – Guangzhou Huali College : Guangzhou North railway station – Lijiao : Longxi – Huangpu Passenger Port : Taihe – Lanhe (→ Foshan) : Xinhe – Jiangnan (→ Dongguan) : Huangpu railway station – Huadu Square Regular-speed metro lines: : Xilang – Guangzhou East railway station : Jiahewanggang – Guangzhou South railway station : Huangcun – Nansha Passenger Port : Jiaokou – Huangpu Passenger Port : Xunfenggang – Guangzhou Middle School : Jiangfu – Haibang : Tanzhonglu – Gaozeng : Gaotangshi – Guanggang New Town (→ Foshan) : Guangzhou railway station – Pazhou – Guangzhou railway station Regular-speed metro lines (continued): : Xunfenggang – Higher Education Mega Center South spur line: Higher Education Mega Center North – Chenbian : Jiaomen – Nansha Passenger Port – Jiaomen : Huangpu railway station – Nanpuxi (→ Foshan) : Lingnan Square – Jiangnan : Chishajiao – Xintang Dadao : Guangzhoudadaobei – Education Park : Dongchong Town – Nansha Wetland Park : Nanguolu – Information Technology Park : Ronggui Railway Station – Qingshengdong : (Foshan →) Huangjinwei – Toubei : Dongjing – Huadong Coach Terminal : Fengcun – Baishantang : Lianxidadao – Shiliuchong : Jiahewanggang – Datian : Shihua – Changping spur line: Yonghe – Lihu : Bicun – Fangshi : Aotou – Conghua Coach Terminal : Nanjiao – GAC Base Foshan : (Foshan →) Guangzhou South Railway Station Foshan : (Foshan →) Xingyedadao Foshan : (Foshan →) Guangzhou Railway Station Foshan : (Foshan →) Baiyun Dongping Foshan : (Foshan →) Longxi Foshan : (Foshan →) Fangcun Foshan : (Foshan →) Hedongdong Dongguan : Huangpu Passenger Port (→ Dongguan) Dongguan : Zengcheng Railway Station (→ Dongguan) Dongguan : Shiqi (→ Dongguan) Connections to neighboring cities The Guangzhou Metro is actively constructing connections to neighboring cities. Foshan is already connected via the Guangfo Metro with connections via Line 7 and Foshan Metro Line 2 is now opened. Dongguan city is proposing connections with Guangzhou Metro Line 13 and the Dongguan Metro. Neighboring Huizhou city proposed in 2016 that Guangzhou Metro Line 16 be extended into Longmen County, achieving the integration of Huizhou and Guangzhou. In January 2018, Huizhou's mayor Mai Jiaomeng revealed that Huizhou was studying two connections with the Guangzhou Metro with Line 16 heading to Yonghan Town, Longmen County and Line 21 extended to Mount Luofu in Boluo County. In 2018, Guangzhou is studying the feasibility of extending Line 18 south into Zhongshan and north into Qingyuan. Guangzhou–Foshan metro connections Fares and tickets Fares Fares of Guangzhou Metro currently range from ¥2 (a couple of stations) to ¥22 (the longest journeys). A journey shorter than 4 km costs ¥2; ¥1 is charged for every 4 km after 4 km, every 6 km after 12 km, and every 8 km after 24 km. Between 30 October 2010 and 30 October 2011, an additional, undiscountable ¥5 fee was charged for any journey to or from Airport South. Collection of such a fee was approved for one year in July 2010 and expired without extension. The fare for the longest possible journey to the exiting station will be charged if a journey exceeds four hours. Passengers may carry luggage below weight and size limits at no cost or a ¥2 surcharge. Current ticket types Single journey ticket Single journey tickets can be bought at a kiosk at every station or at the automatic ticket vending machines. The ticket itself is a contactless radio-frequency plastic token. The user has to tap it on the sensor on the ticket barrier when entering and insert it into a slot at the exit gate where the token is reclaimed. Full base fares are charged for single journey tickets for individuals. Passengers travelling in groups of 30 or larger can enjoy a 10% discount. Yang Cheng Tong and Lingnan Pass Yang Cheng Tong () is a contactless smartcard which can be used on the metro and most other forms of public transport in Guangzhou. Yang Cheng Tong offers discounts for rides on buses and the metro. Within each month, bus and metro rides combined, a 5% discount is available for the first 15 journeys and a 40% discount for all journeys beyond. Full-time students enrolled in primary, secondary, and vocational schools can apply for student passes, which allow them bus and metro rides at half price. Senior citizens can also obtain special passes. Half price is charged for seniors aged 60–64. Seniors aged 65 and above as well as people with major disabilities ride free of charge. Yang Cheng Tong was rebranded in November 2010 as a type of Lingnan Pass (), a new transport card that is valid in multiple cities across the Pearl River Delta. Lingnan Pass cards issued in Guangzhou are named Lingnan Pass·Yang Cheng Tong. Existing cards were automatically upgraded and need not be replaced. Day pass Guangzhou Metro introduced day passes on 1 January 2013. A day pass holder can travel an unlimited number of times in the metro system during a limited period of validity starting from the first use. Two variants are currently available: One-day pass: ¥20 each and valid for 24 hours Three-day pass: ¥50 each and valid for 72 hours Day passes are not rechargeable. They can be fully refunded until the first use, at which time they become nonrefundable. Used passes are not reclaimed, although they can be voluntarily recycled at drop boxes in the stations. The passes are decorated with illustrations of the Cantonese language and cuisine to promote the local culture. The art design was favored by over 70% of those who responded to public opinion surveys compared to two other competing designs. Discontinued ticket types Guangzhou Metro discontinued the following ticket types in favor of Yang Cheng Tong. Stored value ticket Stored value tickets were very similar to Yang Cheng Tong. Stored value tickets are not on sale anymore, but they will be presented as souvenirs to VIPs at the activities of the subway company and can have a 5% discount on fares. Monthly pass Monthly passes were introduced on 1 November 2008 and abolished on 1 May 2010. There were three types of monthly pass: ¥55 monthly pass for 20 single journeys ¥88 monthly pass for 35 single journeys ¥115 monthly pass for 50 single journeys Each journey could travel from one station to any other station regardless of distance. A monthly pass was valid within a calendar month, not the one-month period from the first day it was used. Unused journeys in a month could not be rolled over to a pass for the following month. Student pass and senior citizen pass Both were issued by the metro company and used on metro only, allowing the holders to travel free or at half price. Power supply Most Guangzhou Metro lines in operation are powered by . For power transmission, lines 1, 2, 3, 7, 8, 9 and 13 as well as Guangfo Line use overhead lines, while lines 4, 5, 6, 14 and 21 use third rails. Lines 18 and 22 also use overhead wires, although at . In contrast to the heavy-rail lines, the light-rail APM runs on 600 V 50 Hz 3-phase AC supplied by third rails. Controversies Free rides for relatives of metro employees Starting from 1997 (Guangzhou Metro) implemented a policy that allowed free rides for, in addition to its employees, their relatives. The policy was exposed to the public after its validity was questioned at a hearing on metro fares in December 2005. At first, it was reported that up to three lineal kins of each metro employee were allowed free access to the metro. Based on Guangzhou Metro having about 6,000 employees at the time, participants of the hearing estimated that up to 18,000 relatives of metro employees could ride free at an approximate cost of ¥13 million per year. In response to questions on the policy raised at the hearing, Lu Guanglin, then-General Manager of Guangzhou Metro, claimed that relatives of employees with free access would volunteer as security personnel of the metro. He cited counter-terrorism when explaining that the policy was not exclusively an employee benefit but also a safety measure. Guangzhou Metro later clarified that only the spouse and at most one pre-college child under 18 of each employee were allowed free access, limiting the number of such people to about 2,000. Free rides were strictly regulated and tracked, with abuse subject to disciplinary actions. An unnamed metro employee estimated that the actual cost per year was ¥3 million rather than ¥13 million. Following its publicity, the policy sparked widespread criticism. A Nanfang Daily editorial criticised the policy as Guangzhou Metro exploiting public resources to its own interests. It also questioned the competence of relatives of metro employees in counter-terrorism. It further argued that if Guangzhou Metro indeed needed voluntary security personnel, it could have recruited them openly from the public. Such criticism was echoed by hearing participants as well as members of the Municipal People's Congress of Guangzhou. Guangzhou Metro officially abandoned the policy under pressure on 16 December 2005. Ridership under-prediction The first lines that were constructed, such as Lines 1, 2, and 8, used high capacity 6-car A-type trains in anticipation to heavy ridership. This choice later proved invaluable in the densely populated Guangzhou with all three aforementioned lines today having a peak daily usage of over 1 million passengers each. However, in the early days of operation, ridership of these lines was low. Ridership for Line 1 plateaued at – in the late 1990s and early 2000s even though it was projected to reach in 1998. The under utilization of these lines at the time allowed experts to insist using lower capacity trains on newer lines and even led to the Guangzhou government being criticized for overinflating ridership predictions to approve metro projects. Preference was given small-capacity trains and low-headway operation in the planning of later projects such as Lines 3, 5 and 6. Line 3 was to be built using smaller, lower capacity B-type rolling stock while Lines 5 and 6 was planned to use even lower capacity light metro four car L-type trains. Initially the trains of Line 3 would only be three cars long and planned to gradually be extended into six car trains in the long-term future. This was in line with the conservative ridership projections at the time, with the Airport Section of Line 3 predicted in 2007 to have a long term peak demand of just over 20,000 pphpd by 2034. These ideas would soon prove utterly shortsighted with Line 3 trains being plagued with extreme overcrowding with significant sections of the line over 100% capacity only a few years after opening. Line 3 was forced to adopt its final long term configuration of six-car trains and low headway operation only five years after opening. However, as of 2014, with continuing growth in passenger demand, many sections of Line 3 are still over 100% capacity even after conversion to six car trains and low headway operation. The section crossing the Pearl River between Kecun and Canton Tower stations is the most congested, reaching 136% capacity. In June 2017, the ridership of Line 3 averaged over 2 million passengers per day and on 1 March 2019 the line carried 2.54 million passengers in a single day. With the busiest section carrying over 60,000 pphpd of passenger volume in 2018. As the controversy surrounding Line 3 unfolded the low capacity design of Line 6, another downscaled line, drew concentrated but late criticism from local media in July 2009. Originally believed to have limited attraction to commuters, Line 6 was intended as an auxiliary line with a projected daily ridership of two years after opening and in nine years, These projections assumed the opening year of Line 6 was still 2010 and Guangzhou was less populated. Such projections were in line with ridership of the, at the time, underutilized Lines 1 and Line 2 prior to 2004. However, with the construction of Line 6 well underway using the original plan of four car L-type trains, a change to longer trains had become unrealistic as it would require modification to stations structures whose construction had been completed. An internal report of Guangzhou Metro also released in 2009 reckoned that using the same six car B-type rolling stock as Lines 3 and 7 would increase the capacity of Line 6 by 50%. Land expropriation and residence relocation would pose even greater challenges as evidenced by severe delays in the construction of the stations of Yide Lu and Shahe. In 2014, one year after opening, daily ridership on Line 6 has grown to 600,000 and continues to increase steadily, peaking at 858,000 passengers on 16 September 2016, a mere two years later. With the opening of Phase II extending the line from Changban to Xiangxue in late 2016 ridership continues to increase, averaging 850,000 passengers per day as of April 2018. The congestion following the openings of Lines 3 and 6 made a profound impact on the planning and design of metro lines in Guangzhou. Line 5 had an urgent revision during early construction to support longer six car trains but still using a low capacity L-type design. Lines 7 was originally also planned to use the same four car light metro design as Line 6 but was redesigned and constructed to use higher capacity six car B-type trains. Before the opening of Line 6, the mayor of Guangzhou Chen Jianhua publicly admitted that planning of Line 6 lacked foresight and ridership estimates were too conservative. He predicts the line would be very crowded upon opening. He promised to ensure that future lines will be designed to use trains that are six or more cars long. Newer lines around the city center such as the under construction Line 11, Line 12 and in operation Line 13 will all use high capacity eight car A-type trains. Quality inspection of Line 3 north extension Exposure of quality issue On 11 October 2010, news broke that the concrete structures of two connecting passages in the north extension of Line 3 between Jiahewanggang and Longgui had substandard compressive strength. The quality of the two connecting passages was found to be questionable as early as August 2009. But it not was brought to light until a technician who worked for a company that inspected their quality posted scanned copies of the original inspection reports in his blog in August 2010, and the media picked up the story in October 2010. The connecting passages were intended as connections between two metro tunnels for the maintenance crew and emergency escape corridors for passengers. Their compressive strength was designed to reach 30 MPa. However, the lowest values measured in two inspections were only 21.9 MPa and 25.5 MPa, respectively. Guangzhou Metro and Beijing Chang Cheng Bilfinger Berger Construction Engineering Co., Ltd. (BCBB), contractor of the Jiahewanggang–Longgui section, commissioned two inspection companies to perform a total of three inspections. All three inspections reported results below standard. According to the technician who disclosed the issue and another technician who participated in the first inspection, possible consequences of weaker-than-standard concrete structures included collapse of the passages, blockage of groundwater drains, and even paralysation of the metro tunnels. Alleged fraud attempts According to the two technicians, BCBB rejected a negative inspection report and conspired with their employer company to produce a fraudulent positive report. In response, both the inspection company and BCBB denied their involvement in any fraud attempts. Su Zhenyu, a deputy manager of the Quality and Safety Division of Guangzhou Metro, admitted the quality issue with the connecting passages but maintained the innocence of Guangzhou Metro. According to him (Guangzhou Metro) never received the original inspection reports in 2009 and was unaware of the issue until it received them on 30 September 2010. Su blamed the incident on deceit by BCBB and declared the structures safe for train operation. Su's comments were acknowledged by Guangzhou Metro. Reactions According to Su (Guangzhou Metro) had launched an investigation into the incident and demanded remedial plans for fortifying the structures from the designer after its experts verified that the quality of the passage did not meet the design standard. In its official response (Guangzhou Metro) claimed that it had been monitoring the connecting passages since they were completed in August 2009 and noticed no cracks, deformation or leaks. It also commissioned a re-inspection in September 2010 and obtained results comparable to previous ones. Evaluation by the designer of the connecting passages based on these results recognised their structures as safe. Previously in 2009, the designer also evaluated one of the two connecting passages as safe upon demand of BCBB with the standard for its compressive strength at the lowest permissible value of 25 MPa. In the wake of widespread media coverage, the Construction Commission of Guangzhou launched an investigation into the incident. The commission invited an independent expert group to inspect the connecting passages. The expert group reaffirmed that despite their quality was indeed below the design standard, the passages were safe for operation and needed not be strengthened or rebuilt. The commission also confirmed that BCBB violated regulations in concealing negative inspection reports from related parties. The cause of weaker-than-standard concrete structures was blamed by deputy mayor Su Zequn on cement being mixed manually instead of using machinery due to space limitation at the construction site. The scheduled opening of the north extension of Line 3 on 30 October 2010 was eventually unaffected. Universal free access in November 2010 In January 2010, then-mayor Zhang Guangning revealed to the media that the local government was considering rewarding residents with an "Asian Games gift package" in acknowledgement of their support for the Games. On 27 September 2010, contents of the gift package were officially announced. Included was universal free access to public transit on 30 workdays in November and December 2010 that would coincide with the schedules of the 2010 Asian Games and Asian Para Games in urban areas excluding the districts of Panyu, Nansha and Huadu and the cities of Zengcheng and Conghua. The measure was intended to compensate for the inconvenience caused by a temporary traffic rule that would ban cars from the streets by the parity of the last digits of their license plates during the Games. The free rides policy prompted unprecedented enthusiasm from local residents on 1 November 2010, the first day it went into effect. The metro system carried 7.80 million rides, doubling the figure of an average day. Ridership of the day exceeded the previous peak of 5.13 million on National Day 1 October 2010 by a significant margin and set a national record. Metro traffic remained intense in the days that followed. The daily ridership record was refreshed twice on 3 and 5 November 2010, reaching 7.844 million; total ridership amounted to 38.77 million over the entire workweek. Provisional flow control measures were put into force at all stations, but were utterly inadequate to contain traffic far beyond the design capacity of the metro system. Trains were often crammed, and stations were filled with people queuing in swarms to take a free ride. Guangzhou Metro estimated that when the Asian Games opened, daily ridership would surpass 8 million. Five days after the free rides policy came into force, local authorities decided to rescind the free public transit offer starting from 8 November 2010 and replace it with a cash subsidy program as they deemed the enormous public response a potential security threat to the Games. Registered households and migrant households with presence in the city longer than half a year would each receive a public transit subsidy of ¥150 in cash; individuals in corporate households would each receive ¥50. Residents could claim the subsidies between 12 January and 31 March 2011. Public transit discount policies that were in effect before November 2010 remained unchanged. Kangwang Lu sinkhole incident Around 16:40 on 28 January 2013, in the immediate neighbourhood of the construction site of the Cultural Park Station of Line 6 on Kangwang Lu (), a sinkhole of approximately in area and in depth collapsed, consuming several houses and trees. Six collapses occurred within 40 minutes. Two more collapses occurred later at 21:45, when workers were pouring concrete into the sinkhole. Nearby roads were immediately closed for emergency engineering. The affected section of Kangwang Lu remained closed until the Spring Festival holidays and was closed for a second time on 12 February due to discovery of additional risks. There were no casualties in the incident because metro construction workers detected geological anomalies 20 minutes before the initial collapse and promptly evacuated the neighbourhood. The sinkhole caused disruptions to electricity, gas and water supplies and drainage pipelines. Preliminary analysis blamed the incident on inaccurate geological drawings used for underground blast operations. In total, 412 households, 103 businesses and 69 warehouses were evacuated, and 257 residents were relocated. Guangzhou Metro offered provisional compensations that amounted to ¥50,000 for each collapsed business and ¥2600 for each resident of the collapsed houses, among other compensations. Overseas business On February 25, 2020, the Guangzhou Metro Group and the Punjab Provincial Public Transport Authority of Pakistan signed a service contract for the operation and maintenance of the Orange Line of the Lahore Metro in Pakistan. The bid-winning consortium would undertake the operation and maintenance of the Lahore Metro Orange Line for eight years.	Technology	China	null
1484541	https://en.wikipedia.org/wiki/Wavefront	Wavefront	In physics, the wavefront of a time-varying wave field is the set (locus) of all points having the same phase. The term is generally meaningful only for fields that, at each point, vary sinusoidally in time with a single temporal frequency (otherwise the phase is not well defined). Wavefronts usually move with time. For waves propagating in a unidimensional medium, the wavefronts are usually single points; they are curves in a two dimensional medium, and surfaces in a three-dimensional one. For a sinusoidal plane wave, the wavefronts are planes perpendicular to the direction of propagation, that move in that direction together with the wave. For a sinusoidal spherical wave, the wavefronts are spherical surfaces that expand with it. If the speed of propagation is different at different points of a wavefront, the shape and/or orientation of the wavefronts may change by refraction. In particular, lenses can change the shape of optical wavefronts from planar to spherical, or vice versa. In classical physics, the diffraction phenomenon is described by the Huygens–Fresnel principle that treats each point in a propagating wavefront as a collection of individual spherical wavelets. The characteristic bending pattern is most pronounced when a wave from a coherent source (such as a laser) encounters a slit/aperture that is comparable in size to its wavelength, as shown in the inserted image. This is due to the addition, or interference, of different points on the wavefront (or, equivalently, each wavelet) that travel by paths of different lengths to the registering surface. If there are multiple, closely spaced openings (e.g., a diffraction grating), a complex pattern of varying intensity can result. Simple wavefronts and propagation Optical systems can be described with Maxwell's equations, and linear propagating waves such as sound or electron beams have similar wave equations. However, given the above simplifications, Huygens' principle provides a quick method to predict the propagation of a wavefront through, for example, free space. The construction is as follows: Let every point on the wavefront be considered a new point source. By calculating the total effect from every point source, the resulting field at new points can be computed. Computational algorithms are often based on this approach. Specific cases for simple wavefronts can be computed directly. For example, a spherical wavefront will remain spherical as the energy of the wave is carried away equally in all directions. Such directions of energy flow, which are always perpendicular to the wavefront, are called rays creating multiple wavefronts. The simplest form of a wavefront is the plane wave, where the rays are parallel to one another. The light from this type of wave is referred to as collimated light. The plane wavefront is a good model for a surface-section of a very large spherical wavefront; for instance, sunlight strikes the earth with a spherical wavefront that has a radius of about 150 million kilometers (1 AU). For many purposes, such a wavefront can be considered planar over distances of the diameter of Earth. In an isotropic medium wavefronts travel with the same speed in all directions. Wavefront aberrations Methods using wavefront measurements or predictions can be considered an advanced approach to lens optics, where a single focal distance may not exist due to lens thickness or imperfections. For manufacturing reasons, a perfect lens has a spherical (or toroidal) surface shape though, theoretically, the ideal surface would be aspheric. Shortcomings such as these in an optical system cause what are called optical aberrations. The best-known aberrations include spherical aberration and coma. However, there may be more complex sources of aberrations such as in a large telescope due to spatial variations in the index of refraction of the atmosphere. The deviation of a wavefront in an optical system from a desired perfect planar wavefront is called the wavefront aberration. Wavefront aberrations are usually described as either a sampled image or a collection of two-dimensional polynomial terms. Minimization of these aberrations is considered desirable for many applications in optical systems. Wavefront sensor and reconstruction techniques A wavefront sensor is a device which measures the wavefront aberration in a coherent signal to describe the optical quality or lack thereof in an optical system. There are many applications that include adaptive optics, optical metrology and even the measurement of the aberrations in the eye itself. In this approach, a weak laser source is directed into the eye and the reflection off the retina is sampled and processed. Another application of software reconstruction of the phase is the control of telescopes through the use of adaptive optics. Mathematical techniques like phase imaging or curvature sensing are also capable of providing wavefront estimations. These algorithms compute wavefront images from conventional brightfield images at different focal planes without the need for specialised wavefront optics. While Shack-Hartmann lenslet arrays are limited in lateral resolution to the size of the lenslet array, techniques such as these are only limited by the resolution of digital images used to compute the wavefront measurements. That said, those wavefront sensors suffer from linearity issues and so are much less robust than the original SHWFS, in term of phase measurement. There are several types of wavefront sensors, including: Shack–Hartmann wavefront sensor: a very common method using a Shack–Hartmann lenslet array. Phase-shifting Schlieren technique Wavefront curvature sensor: also called the Roddier test. It yields good correction but needs an already good system as a starting point. Pyramid wavefront sensor Common-path interferometer Foucault knife-edge test Multilateral shearing interferometer Ronchi tester Shearing interferometer Although an amplitude splitting interferometer such as the Michelson interferometer could be called a wavefront sensor, the term is normally applied to instruments that do not require an unaberrated reference beam to interfere with.	Physical sciences	Waves	Physics
1484921	https://en.wikipedia.org/wiki/Xanthidae	Xanthidae	Xanthidae is a family of crabs known as gorilla crabs, mud crabs, pebble crabs or rubble crabs. Xanthid crabs are often brightly coloured and are highly poisonous, containing toxins which are not destroyed by cooking and for which no antidote is known. The toxins are similar to the tetrodotoxin and saxitoxin produced by puffer fish, and may be produced by bacteria in the genus Vibrio living in symbiosis with the crabs, mostly V. alginolyticus and V. parahaemolyticus. Classification Many species formerly included in the family Xanthidae have since been moved to new families. Despite this, Xanthidae is still the largest crab family in terms of species richness, contanining the following subfamilies and genera: Actaeinae Actaea Actaeodes Actaeops † Allactaea Eoxanthops † Epiactaea Epiactaeodes Forestiana Gaillardiellus Heteractaea Lambropsis † Lobiactaea Meractaea Novactaea Odhnea Paractaea Paractaeopsis Phlyctenodes † Platyactaea Pseudoliomera Pseudophlyctenodes † Rata Serenius Banareiinae Banareia Calvactaea Pseudactaea Trichia Chlorodiellinae Chlorodiella Cyclodius Liocarpilodes Luniella Pilodius Ratha Soliella Sulcodius Tweedieia Vellodius Cymoinae Cymo Etisinae Etisus Paraetisus Euxanthinae Alainodaeus Batodaeus Carpoporus Cranaothus Crosnierius Danielea Edwardsium Epistocavea Euxanthus Gothus Guinotellus Hepatoporus Hypocolpus Jacforus Ladomedaeus Lipaesthesius Lipkemedaeus Medaeops Medaeus Miersiella Monodaeus Olenothus Palatigum Paramedaeus Paraxanthodes Pilomedaeus Pleurocolpus Psaumis Rizalthus Takedax Visayax Glyptoxanthinae Glyptoxanthus Kraussiinae Garthasia Kraussia Palapedia Liomerinae Actiomera Bruciana Liomera Lipkemera Neoliomera Neomeria † Paraliomera Polydectinae Lybia Polydectus Tunebia Xanthinae Aldrovandiopanope Aristotelopanope Bottoxanthodes Camilohelleria Cataleptodius Coralliope Cycloxanthops Demania Epixanthops Eurycassiope Euryxanthops Gaudichaudia Guitonia Juxtaxanthias Lachnopodus Leptodius Liagore Lioxanthodes Macromedaeus Marratha Megametope Megamia † Metaxanthops Metopoxantho † Neolioxantho Neoxanthias Neoxanthops Orphnoxanthus Ovatis Palaeoxanthops † Paraxanthias Paraxanthus Pestoxanthodes Pseudomedaeus Wardoxanthops Williamstimpsonia Xanthias Xantho Xanthodius Zosiminae Atergatis Atergatopsis Lophozozymus Paratergatis Platypodia Platypodiella Pulcratis Zosimus Zozymodes Incertae sedis Haydnella † Nogarolia † Sculptoplax †	Biology and health sciences	Crabs and hermit crabs	Animals
1485078	https://en.wikipedia.org/wiki/Desman	Desman	Desmans are aquatic insectivores of the tribe Desmanini (also considered a subfamily, Desmaninae) in the mole family, Talpidae. This tribe consists of two living species found in Europe: the Russian desman (Desmana moschata) in European Russia, and the Pyrenean desman (Galemys pyrenaicus) in the northwest of the Iberian Peninsula and the Pyrenees. Both species are endangered, the Russian desman critically so. They have webbed paws and their front paws are not well-adapted for digging. Desmans were much more diverse and widespread during the Miocene, with two genera, Gaillardia and Magnatalpa, being present in North America. Both living species are thought to have derived from the fossil genus Archaeodesmana. Species Genus Desmana Russian desman (D. moschata) †Desmana kowalskae †Desmana nehringi †Desmana inflata †Desmana thermalis †Desmana marci Genus Galemys Pyrenean desman (G. pyrenaicus) Genus †Asthenoscapter Miocene, Europe Genus †Archaeodesmana Miocene-Pliocene, Europe Genus †Desmanella Miocene, Europe Genus †Gaillardia Miocene, North America Genus †Mygalinia Late Miocene, Hungary Genus †Magnatalpa Miocene-Pliocene, North America Genus †Ruemkelia Gallery In the media	Biology and health sciences	Eulipotyphla	Animals
1486232	https://en.wikipedia.org/wiki/Oarfish	Oarfish	Oarfish are large and extremely long pelagic lampriform fish belonging to the small family Regalecidae. Found in areas spanning from temperate ocean zones to tropical ones, yet rarely seen, the oarfish family contains three species in two genera. One of these, the giant oarfish (Regalecus glesne), is the longest bony fish alive, growing up to about in length. The common name oarfish is thought to allude either to their highly compressed and elongated bodies, or to the now discredited belief that the fish "row" themselves through the water with their pelvic fins. The family name Regalecidae is derived from the Latin regalis, meaning "royal". Although the larger species are considered game fish and are fished commercially to a minor extent, oarfish are rarely caught alive; their flesh is not well regarded for eating due to its gelatinous consistency. Their rarity and large size, and their habit of lingering at the surface when sick or dying, make oarfish a probable source of sea serpent tales. Their beachings after storms have gained them a reputation as harbingers of doom, a folk belief reinforced by the numerous beachings before the disastrous 2011 Tōhoku earthquake and tsunami. Description The dorsal fin originates from above the (relatively large) eyes and runs the entire length of the fish. Of the approximately 400 dorsal fin rays, the first 10 to 13 are elongated to varying degrees, forming a trailing crest embellished with reddish spots and flaps of skin at the ray tips. The pelvic fins are similarly elongated and adorned, reduced to one to five rays each. The pectoral fins are greatly reduced and situated low on the body. The anal fin is completely absent and the caudal fin may be reduced or absent as well, with the body tapering to a fine point. All fins lack true spines. At least one account, from researchers in New Zealand, described the oarfish as giving off "electric shocks" when touched. As in other members of its order, the mouth can be protruded. The body has no scales. In the streamer fish (Agrostichthys parkeri), the skin is clad with hard tubercles; in Regalecus russelii, there are tubercules along the midline of the belly. All species lack swim bladders, the number of gill rakers is variable, but R. russellii has more than R. glesne. Oarfish are silver in coloration; the body is marked with small dark spots. The giant oarfish is by far the largest member of the family, at a length of —with unconfirmed reports of specimens and in length and in weight. The streamer fish reaches in length, while the largest recorded specimen of Regalecus russelii was . Oarfish frequently practise autotomy, self-amputating the tail, presumably as an anti-predator adaptation. All captured R. russellii over 1.5 m long have autotomized tails; it is thought that they may autotomize their tails repeatedly. The break can occur near the tip of the tail so that only a part of the caudal fin is lost, or it may involve a few caudal vertebrae; in extreme cases the entire tail is lost. The wound heals but the tail does not regenerate. Hyperossified bones have been documented in several oarfish washed up on the coast of California. These bony rays run along the entire dorsal length of the body. Their function is both to provide structural support to the spine during undulations (tail movement used for locomotion), and to prevent stress fractures that could occur from strong movement. Unlike many deep-sea fish, oarfish have no swim bladders for maintaining depth in the water column. It is likely that this forces more frequent tail undulations as the main mode of depth regulation in oarfish. Evolution Phylogeny Through the analysis of the mitochondrial genome of Regalecus glesne, the phylogenetic placement of the giant oarfish was further verified. Oarfish are Lampriformes, so placed due to their morphology. Analysis of the mitochondrial genome of an R. glesne specimen clusters the species with Trachipterus trachypterus and Zu cristatus, two other Lampriformes. Taxonomy Oarfish were first described in 1772. Three extant species in two extant genera are described: Giant Oarfish (Regalecus glesne) Russell's Oarfish (Regalecus russelii) Streamerfish (Agrostichthys parkeri) Environment and distribution The oarfish inhabits the epipelagic to mesopelagic ocean layers, ranging from 250 meters (660 ft) to 1,000 meters (3,300 ft) and is rarely seen on the surface. A few have been found still barely alive, but usually if one floats to the surface, it dies due to depressurisation. At the depths the oarfish live, there are few or no currents. As a result, they build little muscle mass and they cannot survive in shallower turbulent water. The members of the family have a worldwide range, with tropical, subtropical, and warm temperate distributions. The oarfish typically reside in the mesopelagic area of the sea. However, human encounters with live oarfish are rare, and distribution information is collated from records of oarfish caught or washed ashore. Ecology and life history Behaviour Rare encounters with divers and accidental catches have supplied what little is known of oarfish ethology (behavior) and ecology. In 2001, an oarfish was filmed alive in the wild. The fish was spotted by a group of U.S. Navy personnel during the inspection of a buoy in the Bahamas. The oarfish was observed to propel itself by an amiiform mode of swimming; that is, rhythmically undulating the dorsal fin while keeping the body itself straight. Perhaps indicating a feeding posture, oarfish have been observed swimming in a vertical orientation. In this posture, the downstreaming light would silhouette the oarfishes' prey, making them easier to spot. An oarfish measuring and was caught in February 2003 using a fishing rod baited with squid at Skinningrove, United Kingdom. In July 2008, scientists for the first time captured footage of an oarfish swimming in its natural habitat in the mesopelagic zone in the Gulf of Mexico. The fish was estimated to be between in length. Five observations of apparently healthy oarfish Regalecus glesne by remotely operated vehicles were reported from the northern Gulf of Mexico between 2008 and 2011 at depths within the epipelagic and mesopelagic zones. These observations include the deepest verified record of R. glesne (). In the 2011 sighting, an oarfish has been observed to switch from swimming with a vertical posture to swimming laterally, using lateral undulations of its entire body. Oarfish were found to have late or slow flight responses towards approaching remotely operated vehicles, supporting the hypothesis that they have few natural predators. From December 2009 to March 2010, unusual numbers of the slender oarfish Regalecus russelii appeared in the waters and on the beaches of Japan. In 2016, Animal Planet aired an episode of the television series River Monsters named "Deep Sea Demon" in which Jeremy Wade was filmed with a live oarfish. The oarfish at this location seemed to be using a buoy anchor chain as a guide to ascend to the surface. On his second diving attempt, he filmed two live oarfish as they came relatively close to the surface. Wade was able to touch one of the oarfish with his hand. In January 2019 two oarfish were found alive in the nets of fishermen on the Japanese island of Okinawa. Feeding ecology Oarfish feed primarily on zooplankton, selectively straining tiny euphausiids, shrimp, and other crustaceans from the water. Small fish, jellyfish, and squid are also taken. It has been observed that oarfish eat by suctioning prey such as plankton blooms while in the water. Reproduction and life history The oceanodromous Regalecus glesne is recorded as spawning off Mexico from July to December; all species are presumed to not guard their eggs, and release brightly coloured, buoyant eggs, up to across, which are incorporated into the zooplankton. Based on their reproductive morphology, oarfish are thought to batch spawn. Within each breeding season that may last one or two months, individuals spawn once or multiple times in discrete spawning events before their gonads enter a long, regressive stage of reproductive development. The eggs hatch after about three weeks into highly active larvae that feed on other zooplankton. The larvae have little resemblance to the adults, with long dorsal and pelvic fins and extensible mouths. Larvae and juveniles have been observed drifting just below the surface. In contrast, adult oarfish are rarely seen at the surface when not sick or injured. It is probable that the fishes go deeper as they mature. From January to February 2019, researchers tested and recorded the first successful instance of artificial insemination and hatching of Regalecus russellii using gonads from two washed-up specimens. Compared to adults, the body structure of newly hatched oarfish larvae look more compressed. The larvae often swam using mainly their pectoral fins, facing downward, with their mouths constantly open. The larvae were invertebrates but had bones in their head area, as well as fins. They died of starvation four days after they hatched. Female R. russelii have bifurcated ovaries with a cavity through which the eggs pass before they are laid. The testes of male oarfish in the coelomic cavity near the digestive tract. There are two separate, disconnected testes, the left one being longer than the one on the right. A single female can produce hundreds of thousands, to millions of eggs. The eggs are laid in the water column, and they float freely in the water. Predators and parasites A 2015 study suggested that the shortfin mako shark and the sperm whale could both be predators of the oarfish, based on patterns of parasite transmission and analysis of oarfish viscera. In folklore The slender oarfish, (竜宮の使い "Ryūgū-No-Tsukai"), known in Japanese folklore as the 'Messenger from the Sea God's Palace', is said to portend earthquakes. The oarfish has been nicknamed the "doomsday fish" because, historically, appearances of the fish were linked with subsequent natural disasters, namely earthquakes or tsunamis. After the 2011 Tōhoku earthquake and tsunami which killed over 20,000 people, many in Japan pointed to the 20 oarfish washed up on the country's beaches in 2009 and 2010 in line with this reputation as a harbinger of doom.	Biology and health sciences	Fishes	null
1487479	https://en.wikipedia.org/wiki/Lobster%20trap	Lobster trap	A lobster trap or lobster pot is a portable trap that traps lobsters or crayfish and is used in lobster fishing. In Scotland (chiefly in the north), the word creel was used to refer to a device used to catch lobsters and other crustaceans. A lobster trap can hold several lobsters. Lobster traps can be constructed of wire and wood, metal and netting, or rigid plastic. An opening permits the lobster to enter a tunnel of netting or other one-way device. Pots are sometimes constructed in two parts, called the "chamber" or "kitchen", where there is bait, and exits into the "parlor", which prevents escape. Lobster pots are usually dropped to the sea floor, one or more at a time, sometimes up to 40 or more, and are marked by a buoy so they can be picked up later. Description The trap can consist of a wood frame surrounded by mesh. The majority of the newer traps found in the Northeast of the US and the Canadian Maritimes consist of a plastic-coated metal frame. A piece of bait, often fish or chum, is placed inside the trap, and the traps are dropped onto the sea floor. A long rope is attached to each trap, at the end of which is a plastic or styrofoam buoy that bears the owner's license number. The entrances to the traps are designed to be one-way entrances only. The traps are checked every other day by the fisherman and rebaited if necessary. One study indicated that lobster traps are very inefficient and allow almost all lobsters to escape. Automatic rebaiting improves efficiency. History The lobster trap was invented in 1808 by Ebenezer Thorndike of Swampscott, Massachusetts. By 1810, the wooden lath trap is said to have originated in Cape Cod, Massachusetts. New England fishermen in the United States used it for years before American companies introduced it to the Canadian fishery through their Atlantic coast canneries. An 1899 report by the United States Fish Commission on the Lobster Fishery Of Maine, described the local "lath pots" used by Maine lobster fishers: Safety Lobster fishermen who become entangled in the trap line are at risk of drowning if they are pulled overboard. Best practices have been developed to prevent and reduce entanglement and to facilitate getting fishermen who have fallen overboard back onto their vessels. Rope-less lobster traps As whales can get entangled in ropes, there is currently research going on into the development of rope-less lobster traps. Some designs have already been developed.	Technology	Hunting and fishing	null
1487891	https://en.wikipedia.org/wiki/GDK	GDK	GDK (GIMP Drawing Kit) is a library that acts as a wrapper around the low-level functions provided by the underlying windowing and graphics systems. GDK lies between the display server and the GTK library, handling basic rendering such as drawing primitives, raster graphics (bitmaps), cursors, fonts, as well as window events and drag-and-drop functionality. Like GTK Scene Graph Kit (GSK), GDK is part of GTK and licensed under the GNU Lesser General Public License (LGPL). Software architecture GTK is implemented on top of an abstraction layer called GDK, freeing GTK from low-level concerns like input gathering, Drag and drop and pixel format conversion. GDK is an intermediate layer which separates GTK from the details of the windowing system. GDK is an important part of GTK's portability. Since low-level cross-platform functionality is already provided by GLib, all that is needed to make GTK run on other platforms is to port GDK to the underlying operating system's graphics layer. Hence, the GDK ports to the Windows API and Quartz are what enable GTK applications to run on Windows and macOS, respectively. Starting with GTK+ 2.8, GDK supports Cairo which should be used with GTK+ 3 instead of GDK's drawing functions. GDK is an intermediate layer which isolates GTK from the details of the windowing system. GDK is a thin wrapper around Xlib. The X Window System comes with a low-level library called Xlib. Almost every function in GDK is a very thin wrapper around a corresponding Xlib function; but some of the complexity (and functionality) of Xlib is hidden, to simplify programming and to make GDK easier to port to other windowing systems, such as Wayland or Microsoft Windows. The concealed Xlib functionality will rarely be of interest to application programmers; for example, many features used solely by window managers are not exposed in GDK. GDK lets you do low level stuff, like e.g. "blit this pixmap to the screen". GDK provides a layer that is much more portable than say the X protocol, without sacrificing any of the low-level accessibility that systems such as X provide. The true power of this abstraction is that if you choose to use it rather than say, X, your software will automatically render on the Linux Framebuffer and Windows. Having OpenGL (or OpenGL ES) support in GDK, facilitates a slightly better control of the graphics pipeline; OpenGL is well suited for compositing textured data but totally unsuited for drawing. GdkFrameClock GdkFrameClock was added in GTK 3.8 While GTK applications remain mainloop driven (cf. Glib event loop), meaning the application is idle inside this main loop most of the time and just waits for something to happen and then calls the appropriate subroutine when it does, GdkFrameClock adds an additional mechanism, that gives a "pulse" to the application. It tells the application when to update and repaint a window. The beat rate can be synchronized with the monitor refresh rate. GTK Scene Graph Kit In its history GDK contained and linked with a couple of different Canvases. https://wiki.gnome.org/Attic/ProjectRidley/CanvasOverview https://wiki.gnome.org/Attic/ProjectRidley/CanvasOverview/Canvases https://wiki.gnome.org/Projects/GooCanvas Developers were also considering new directions for the library, including removing deprecated API components and adding an integrated scene graph (canvas) system, similar to the Clutter graphics library, effectively integrating GTK with OpenGL and Vulkan. GTK Scene Graph Kit (GSK) GTK+ Scene Graph Kit (GSK) was released as part of GTK+ 3.90 in March 2017. It is the scene graph and rendering API for GTK. GSK has not been further integrated with GDK (which is also part of GTK) but is kept in its own directory. Windowing systems GDK contains back-ends to a couple of windowing systems, namely to the X11 and Wayland protocols, to Quartz and GDI, and even to the Hypertext Transfer Protocol (HTTP) engine Broadway. With the release of GNOME 3.16 in March 2015, GDK obtained an experimental back-end for the Mir display server protocol. The Mir display server protocol is a product by Canonical for their Ubuntu distribution of Linux, which they intend to compete with the Wayland display server protocol; so far, it is implemented only in Ubuntu. At present, no back-end exists for KMS. To start an application and force this instance of it to use a certain windowing system, you specify the variable GDK_BACKEND: GDK_BACKEND=wayland gnome-calculator GDK_BACKEND=wayland CLUTTER_BACKEND=wayland cheese gdk-pixbuf gdk-pixbuf is a toolkit for image loading and pixel buffer manipulation. The library provides image loading and saving facilities, fast scaling and compositing of pixbufs, simple animation loading (i.e. animated GIFs), and rendering the libart image buffer to a GdkDrawable instance. gdk-pixbuf has a fairly large API. The fundamental structure in the gdk-pixbuf library is GdkPixbuf, a private, opaque data structure that mirrors many of the same concepts that ArtPixBuf supports. In fact, most of GdkPixbuf's private data fields have the same names and data types as the corresponding ones in ArtPixBuf. This similarity dates back to the earlier days when gdk-pixbuf was a wrapper around libart. Since that time, the libart dependency has been stripped out, and gdk-pixbuf was merged into the GTK+ 2.0 code base. As such, gdk-pixbuf's days as a standalone library are limited to the GNOME 1 release. With the release of GTK+ 2.22 on 2010-09-23, gdk-pixbuf has been turned back into a standalone library, after being shipped as part of GTK+ since gtk+ 2.0. This was done in preparation for the transition to GTK+ 3. https://git.gnome.org/browse/gdk-pixbuf/ The first stand-alone release was 2.22 on 2010-Sep-21, its development started with 2.21.3 on 2010-06-23. History GDK was originally developed on the X Window System for the GIMP raster graphics editor.	Technology	System	null
1488243	https://en.wikipedia.org/wiki/Sodium%20acetate	Sodium acetate	Sodium acetate, CH3COONa, also abbreviated NaOAc, is the sodium salt of acetic acid. This salt is colorless deliquescent, and hygroscopic. Applications Biotechnological Sodium acetate is used as the carbon source for culturing bacteria. Sodium acetate can also be useful for increasing yields of DNA isolation by ethanol precipitation. Industrial Sodium acetate is used in the textile industry to neutralize sulfuric acid waste streams and also as a photoresist while using aniline dyes. It is also a pickling agent in chrome tanning and helps to impede vulcanization of chloroprene in synthetic rubber production. It is also used to reduce static electricity during production of disposable cotton pads. Concrete longevity Sodium acetate is used to mitigate water damage to concrete by acting as a concrete sealant, while also being environmentally benign and cheaper than the commonly used epoxy alternative for sealing concrete against water permeation. Food Sodium acetate (anhydrous) is widely used as a shelf-life extending agent and pH-control agent. It is safe to eat at low concentration. Buffer solution A solution of sodium acetate (a basic salt of acetic acid) and acetic acid can act as a buffer to keep a relatively constant pH level. This is useful especially in biochemical applications where reactions are pH-dependent in a mildly acidic range (pH 4–6). Heating pad Sodium acetate is also used in heating pads, hand warmers, and hot ice. A supersaturated solution of sodium acetate in water is supplied with a device to initiate crystallization, a process that releases substantial heat. Sodium acetate trihydrate crystals melt at , and the liquid sodium acetate dissolves in the released water of crystallization. When heated past the melting point and subsequently allowed to cool, the aqueous solution becomes supersaturated. This solution is capable of cooling to room temperature without forming crystals. By pressing on a metal disc within the heating pad, a nucleation center is formed, causing the solution to crystallize back into solid sodium acetate trihydrate. The process of crystallization is exothermic. The latent heat of fusion is about 264–289 kJ/kg. Unlike some types of heat packs, such as those dependent upon irreversible chemical reactions, a sodium acetate heat pack can be easily reused by immersing the pack in boiling water for a few minutes, until the crystals are completely dissolved, and allowing the pack to slowly cool to room temperature. Preparation For laboratory use, sodium acetate is inexpensive and usually purchased instead of being synthesized. It is sometimes produced in a laboratory experiment by the reaction of acetic acid, commonly in the 5–18% solution known as vinegar, with sodium carbonate ("washing soda"), sodium bicarbonate ("baking soda"), or sodium hydroxide ("lye", or "caustic soda"). Any of these reactions produce sodium acetate and water. When a sodium and carbonate ion-containing compound is used as the reactant, the carbonate anion from sodium bicarbonate or carbonate, reacts with the hydrogen from the carboxyl group (-COOH) in acetic acid, forming carbonic acid. Carbonic acid readily decomposes under normal conditions into gaseous carbon dioxide and water. This is the reaction taking place in the well-known "volcano" that occurs when the household products, baking soda and vinegar, are combined. CH3COOH + NaHCO3 → CH3COONa + → + Industrially, sodium acetate trihydrate is prepared by reacting acetic acid with sodium hydroxide using water as the solvent. CH3COOH + NaOH → CH3COONa + H2O. To manufacture anhydrous sodium acetate industrially, the Niacet Process is used. Sodium metal ingots are extruded through a die to form a ribbon of sodium metal, usually under an inert gas atmosphere such as N2 then immersed in anhydrous acetic acid. 2 CH3COOH + 2 Na →2 CH3COONa + H2. The hydrogen gas is normally a valuable byproduct. Structure The crystal structure of anhydrous sodium acetate has been described as alternating sodium-carboxylate and methyl group layers. Sodium acetate trihydrate's structure consists of distorted octahedral coordination at sodium. Adjacent octahedra share edges to form one-dimensional chains. Hydrogen bonding in two dimensions between acetate ions and water of hydration links the chains into a three-dimensional network. Reactions Sodium acetate can be used to form an ester with an alkyl halide such as bromoethane: CH3COONa + BrCH2CH3 → CH3COOCH2CH3 + NaBr Sodium acetate undergoes decarboxylation to form methane (CH4) under forcing conditions (pyrolysis in the presence of sodium hydroxide): CH3COONa + NaOH → CH4 + Na2CO3 Calcium oxide is the typical catalyst used for this reaction. Cesium salts also catalyze this reaction.	Physical sciences	Acetates	Chemistry
1488463	https://en.wikipedia.org/wiki/Agroforestry	Agroforestry	Agroforestry (also known as agro-sylviculture or forest farming) is a land use management system that integrates trees with crops or pasture. It combines agricultural and forestry technologies. As a polyculture system, an agroforestry system can produce timber and wood products, fruits, nuts, other edible plant products, edible mushrooms, medicinal plants, ornamental plants, animals and animal products, and other products from both domesticated and wild species. Agroforestry can be practiced for economic, environmental, and social benefits, and can be part of sustainable agriculture. Apart from production, benefits from agroforestry include improved farm productivity, healthier environments, reduction of risk for farmers, beauty and aesthetics, increased farm profits, reduced soil erosion, creating wildlife habitat, less pollution, managing animal waste, increased biodiversity, improved soil structure, and carbon sequestration. Agroforestry practices are especially prevalent in the tropics, especially in subsistence smallholdings areas, with particular importance in sub-Saharan Africa. Due to its multiple benefits, for instance in nutrient cycle benefits and potential for mitigating droughts, it has been adopted in the USA and Europe. Definition At its most basic, agroforestry is any of various polyculture systems that intentionally integrate trees with crops or pasture on the same land. An agroforestry system is intensively managed to optimize helpful interactions between the plants and animals included, and “uses the forest as a model for design." Agroforestry shares principles with polyculture practices such as intercropping, but can also involve much more complex multi-strata agroforests containing hundreds of species. Agroforestry can also utilise nitrogen-fixing plants such as legumes to restore soil nitrogen fertility. The nitrogen-fixing plants can be planted either sequentially or simultaneously. History and scientific study The term “agroforestry” was coined in 1973 by Canadian forester John Bene, but the concept includes agricultural practices that have existed for millennia. Scientific agroforestry began in the 20th century with ethnobotanical studies carried out by anthropologists. However, indigenous communities that have lived in close relationships with forest ecosystems have practiced agroforestry informally for centuries. For example, Indigenous peoples of California periodically burned oak and other habitats to maintain a ‘pyrodiversity collecting model,’ which allowed for improved tree health and habitat conditions. Likewise Native Americans in the eastern United States extensively altered their environment and managed land as a “mosaic” of woodland areas, orchards, and forest gardens. Agroforestry in the tropics is ancient and widespread throughout various tropical areas of the world, notably in the form of "tropical home gardens." Some “tropical home garden” plots have been continuously cultivated for centuries. A “home garden” in Central America could contain 25 different species of trees and food crops on just one-tenth of an acre. "Tropical home gardens" are traditional systems developed over time by growers without formalized research or institutional support, and are characterized by a high complexity and diversity of useful plants, with a canopy of tree and palm species that produce food, fuel, and shade, a mid-story of shrubs for fruit or spices, and an understory of root vegetables, medicinal herbs, beans, ornamental plants, and other non-woody crops. In 1929, J. Russel Smith published Tree Crops: A Permanent Agriculture, in which he argued that American agriculture should be changed two ways: by using non-arable land for tree agriculture, and by using tree-produced crops to replace the grain inputs in the diets of livestock. Smith wrote that the honey locust tree, a legume that produced pods that could be used as nutritious livestock feed, had great potential as a crop. The book's subtitle later led to the coining of the term permaculture. The most studied agroforestry practices involve a simple interaction between two components, such as simple configurations of hedges or trees integrated with a single crop. There is significant variation in agroforestry systems and the benefits they have. Agroforestry as understood by modern science is derived from traditional indigenous and local practices, developed by living in close association with ecosystems for many generations. Benefits Benefits include increasing farm productivity and profitability, reduced soil erosion, creating wildlife habitat, managing animal waste, increased biodiversity, improved soil structure, and carbon sequestration. Agroforestry systems can provide advantages over conventional agricultural and forest production methods. They can offer increased productivity; social, economic and environmental benefits, as well as greater diversity in the ecological goods and services provided. These benefits are conditional on good farm management. This includes choosing the right trees, as well as pruning them regularly etc. Biodiversity Biodiversity in agroforestry systems is typically higher than in conventional agricultural systems. Two or more interacting plant species in a given area create a more complex habitat supporting a wider variety of fauna. Agroforestry is important for biodiversity for different reasons. It provides a more diverse habitat than a conventional agricultural system in which the tree component creates ecological niches for a wide range of organisms both above and below ground. The life cycles and food chains associated with this diversification initiate an agroecological succession that creates functional agroecosystems that confer sustainability. Tropical bat and bird diversity, for instance, can be comparable to the diversity in natural forests. Although agroforestry systems do not provide as many floristic species as forests and do not show the same canopy height, they do provide food and nesting possibilities. A further contribution to biodiversity is that the germplasm of sensitive species can be preserved. As agroforests have no natural clear areas, habitats are more uniform. Furthermore, agroforests can serve as corridors between habitats. Agroforestry can help conserve biodiversity, positively influencing other ecosystem services. Soil and plant growth Depleted soil can be protected from soil erosion by groundcover plants such as naturally growing grasses in agroforestry systems. These help to stabilise the soil as they increase cover compared to short-cycle cropping systems. Soil cover is a crucial factor in preventing erosion. Cleaner water through reduced nutrient and soil surface runoff can be a further advantage of agroforestry. Trees can help reduce water runoff by decreasing water flow and evaporation and thereby allowing for increased soil infiltration. Compared to row-cropped fields nutrient uptake can be higher and reduce nutrient loss into streams. Further advantages concerning plant growth: Bioremediation Drought tolerance Increased crop stability Sustainability Agroforestry systems can provide ecosystem services which can contribute to sustainable agriculture in the following ways: Diversification of agricultural products, such as fuelwood, medicinal plants, and multiple crops, increases income security Increased food security and nutrition by restored soil fertility, crop diversity and resilience to weather shocks for food crops Land restoration through reducing soil erosion and regulating water availability Multifunctional site use, e.g., crop production and animal grazing Reduced deforestation and pressure on woodlands by providing farm-grown fuelwood Possibility of reduced chemicals inputs, e.g. due to improved use of fertilizer, increased resilience against pests, and increased ground cover which reduces weeds Growing space for medicinal plants e.g., in situations where people have limited access to mainstream medicines According to the United Nations Food and Agriculture Organization (FAO)'s The State of the World’s Forests 2020, adopting agroforestry and sustainable production practices, restoring the productivity of degraded agricultural lands, embracing healthier diets and reducing food loss and waste are all actions that urgently need to be scaled up. Agribusinesses must meet their commitments to deforestation-free commodity chains and companies that have not made zero-deforestation commitments should do so. Other environmental goals Carbon sequestration is an important ecosystem service. Agroforestry practices can increase carbon stocks in soil and woody biomass. Trees in agroforestry systems, like in new forests, can recapture some of the carbon that was lost by cutting existing forests. They also provide additional food and products. The rotation age and the use of the resulting products are important factors controlling the amount of carbon sequestered. Agroforests can reduce pressure on primary forests by providing forest products. Adaptation to climate change Agroforestry can significantly contribute to climate change mitigation along with adaptation benefits. A case study in Kenya found that the adoption of agroforestry drove carbon storage and increased livelihoods simultaneously among small-scale farmers. In this case, maintaining the diversity of tree species, especially land use and farm size are important factors. Poor smallholder farmers have turned to agroforestry as a means to adapt to climate change. A study from the CGIAR research program on Climate Change, Agriculture and Food Security found from a survey of over 700 households in East Africa that at least 50% of those households had begun planting trees in a change from earlier practices. The trees were planted with fruit, tea, coffee, oil, fodder and medicinal products in addition to their usual harvest. Agroforestry was one of the most widespread adaptation strategies, along with the use of improved crop varieties and intercropping. Tropical Trees in agroforestry systems can produce wood, fruits, nuts, and other useful products. Agroforestry practices are most prevalent in the tropics, especially in subsistence smallholdings areas such as sub-Saharan Africa. Research with the leguminous tree Faidherbia albida in Zambia showed maximum maize yields of 4.0 tonnes per hectare using fertilizer and inter-cropped with the trees at densities of 25 to 100 trees per hectare, compared to average maize yields in Zimbabwe of 1.1 tonnes per hectare. Hillside systems A well-studied example of an agroforestry hillside system is the Quesungual Slash and Mulch Agroforestry System in Lempira Department, Honduras. This region was historically used for slash-and-burn subsistence agriculture. Due to heavy seasonal floods, the exposed soil was washed away, leaving infertile barren soil exposed to the dry season. Farmed hillside sites had to be abandoned after a few years and new forest was burned. The UN's FAO helped introduce a system incorporating local knowledge consisting of the following steps: Thin and prune Hillside secondary forest, leaving individual beneficial trees, especially nitrogen-fixing trees. They help reduce soil erosion, maintain soil moisture, provide shade and provide an input of nitrogen-rich organic matter in the form of litter. Plant maize in rows. This is a traditional local crop. Harvest from the dried plant and plant beans. The maize stalks provide an ideal structure for the climbing bean plants. Bean is a nitrogen-fixing plant and therefore helps introduce more nitrogen. Pumpkins can be planted during this time. The plant's large leaves and horizontal growth provide additional shade and moisture retention. It does not compete with the beans for sunlight since the latter grow vertically on the stalks. Every few seasons, rotate the crop by grazing cattle, allowing grass to grow and adding soil organic matter and nutrients (manure). The cattle prevent total reforestation by grazing around the trees. Repeat. Kuojtakiloyan The kuojtakiloyan of Mexico is a jungle-landscaped polyculture that grows avocadoes, sweet potatoes, cinnamon, black cherries, , citrus fruits, gourds, macadamia, mangoes, bananas and sapotes. Kuojtakiloyan is a Masehual term that means 'useful forest' or 'forest that produces', and it is an agroforestry system developed and maintained by indigenous peoples of the Sierra Norte of the State of Puebla, Mexico. It has become a vital fountain of resources (food, medicinal herbs, fuels, floriculture, etc.) for the local population, but it is also a respectful transformation of the environment, with its biodiversity and nature conservation. The kuojtakiloyan comes directly from the ancestral Nahua and Totonaku knowledge of their natural environment. Despite its unawareness among the mainstream Mexican population, many agronomic experts in the world point it out as a successful case of sustainable agroforestry practiced communally. The kuojtakiloyan is a jungle-landscaped polyculture in which avocados, sweet potatoes, cinnamon, black cherries, chalahuits, citrus fruits, gourds, macadamia, mangoes, bananas and sapotes are grown. In addition, a wide variety of harvested wild edible mushrooms and herbs (quelites). The jonote is planted because its fiber is useful in basketry, and also bamboo, which is fast growing, to build cabins and other structures. Concurrently to kuojtakiloyan, shade coffee is grown (café bajo sombra in Spanish; kafentaj in Masehual). Shade is essential to obtain high quality coffee. The local population has favored the proliferation of the stingless bee (pisilnekemej) by including the plants that it pollinates. From bees, they get honey, pollen, wax and propolis. Shade crops With shade applications, crops are purposely raised under tree canopies within the shady environment. The understory crops are shade tolerant or the overstory trees have fairly open canopies. A conspicuous example is shade-grown coffee. This practice reduces weeding costs and improves coffee quality and taste. Crop-over-tree systems Crop-over-tree systems employ woody perennials in the role of a cover crop. For this, small shrubs or trees pruned to near ground level are utilized. The purpose is to increase in-soil nutrients and/or to reduce soil erosion. Intercropping and alley cropping With alley cropping, crop strips alternate with rows of closely spaced tree or hedge species. Normally, the trees are pruned before planting the crop. The cut leafy material - for example, from Alchornea cordifolia and Acioa barteri - is spread over the crop area to provide nutrients. In addition to nutrients, the hedges serve as windbreaks and reduce erosion. In tropical areas of North and South America, various species of Inga such as I. edulis and I. oerstediana have been used for alley cropping. Intercropping is advantageous in Africa, particularly in relation to improving maize yields in the sub-Saharan region. Use relies upon the nitrogen-fixing tree species Sesbania sesban, Tephrosia vogelii, Gliricidia sepium and Faidherbia albida. In one example, a ten-year experiment in Malawi showed that, by using the fertilizer tree Gliricidia (G. sepium) on land on which no mineral fertilizer was applied, maize/corn yields averaged as compared to in plots without fertilizer trees or mineral fertilizers. Weed control is inherent to alley cropping, by providing mulch and shade. Syntropic systems Syntropic farming, syntropic agriculture or syntropic agroforestry is an organic, permaculture agroforestry system developed by Ernst Götsch in Brazil. Sometimes this system is referred to as a successional agroforestry systems or SAFS, which sometimes refer to a broader concept originating in Latin America. The system focuses on replicating natural systems of accumulation of nutrients in ecosystems, replicating secondary succession, in order to create productive forest ecosystems that produce food, ecosystem services and other forest products. The system relies heavily on several processes: Dense planting mixing perennial and annual crops Rapid cutting and composting of fast growing pioneer species, to accumulate nutrients and biomass Creating greater water retention on the land through improving penetration of water into soil and plant water cycling The systems were first developed in tropical Brazil, but many similar systems have been tested in temperate environments as soil and ecosystem restoration tactics. The framework for the syntropic agroforestry is advocated for by Agenda Gotsch an organization built to promote the systems. Syntropic systems have a number of documented benefits, including increased soil water penetration, increases to productivity on marginal land that has since become and soil temperature moderation. In Burma Taungya is a system from Burma. In the initial stages of an orchard or tree plantation, trees are small and widely spaced. The free space between the newly planted trees accommodates a seasonal crop. Instead of costly weeding, the underutilized area provides an additional output and income. More complex taungyas use between-tree space for multiple crops. The crops become more shade tolerant as the tree canopies grow and the amount of sunlight reaching the ground declines. Thinning can maintain sunlight levels. In India Itteri agroforestry systems have been used in Tamil Nadu since time immemorial. They involve the deliberate management of multipurpose trees and shrubs grown in intimate association with herbaceous species. They are often found along village and farm roads, small gullies, and field boundaries. Bamboo-based agroforestry systems (Dendrocalamus strictus + sesame–chickpea) have been studied for enhancing productivity in semi-arid tropics of central India. In Africa A project to mitigate climate change with agriculture was launched in 2019 by the "Global EverGreening Alliance". The target is to sequester carbon from the atmosphere. By 2050 the restored land should sequestrate 20 billion tons of carbon annually Shamba (Swahili for 'plantation') is an agroforestry system practiced in East Africa, particularly in Kenya. Under this system, various crops are combined: bananas, beans, yams and corn, to which are added timber resources, beekeeping, medicinal herbs, mushrooms, forest fruits, fodder for livestock, etc. In Hawai'i Native Hawaiians formerly practiced agroforestry adapted to the islands' tropical landscape. Their ability to do this influenced the region's carrying capacity, social conflict, cooperation, and political complexity. More recently, after scientific study of lo’I systems, attempts have been made to reintroduce dryland agroforestry in Hawai’i Island and Maui, fostering interdisciplinary collaboration between political leaders, landowners, and scientists. Temperate Although originally a concept in tropical agronomy, agroforestry's multiple benefits, for instance in nutrient cycles and potential for mitigating droughts, have led to its adoption in the USA and Europe. The United States Department of Agriculture distinguishes five applications of agroforestry for temperate climates, namely alley cropping, forest farming, riparian forest buffers, silvopasture, and windbreaks. Alley cropping Alley cropping can also be used in temperate climates. Strip cropping is similar to alley cropping in that trees alternate with crops. The difference is that, with alley cropping, the trees are in single rows. With strip cropping, the trees or shrubs are planted in wide strips. The purpose can be, as with alley cropping, to provide nutrients, in leaf form, to the crop. With strip cropping, the trees can have a purely productive role, providing fruits, nuts, etc. while, at the same time, protecting nearby crops from soil erosion and harmful winds. Inga alley cropping Inga alley cropping is the planting agricultural crops between rows of Inga trees. It has been promoted by Mike Hands. Using the Inga tree for alley cropping has been proposed as an alternative to the much more ecologically destructive slash and burn cultivation. The technique has been found to increase yields. It is sustainable agriculture as it allows the same plot to be cultivated over and over again thus eliminating the need for burning of the rainforests to get fertile plots. Inga tree Inga trees are native to many parts of Central and South America. Inga grows well on the acid soils of the tropical rainforest and former rainforest. They are leguminous and fix nitrogen into a form usable by plants. Mycorrhiza growing within the roots (arbuscular mycorrhiza) was found to take up spare phosphorus, allowing it to be recycled into the soil. Other benefits of Inga include the fact that it is fast growing with thick leaves which, when left on the ground after pruning, form a thick cover that protects both soil and roots from the sun and heavy rain. It branches out to form a thick canopy so as to cut off light from the weeds below and withstands careful pruning year after year. History The technique was first developed and trialled by tropical ecologist Mike Hands in Costa Rica in the late 1980s and early '90s. Research funding from the EEC allowed him to experiment with species of Inga. Although alley cropping had been widely researched, it was thought that the tough pinnate leaves of the Inga tree would not decompose quickly enough. The Inga is used as hedges and pruned when large enough to provide a mulch in which bean and corn seeds are planted. This results in both improving crop yields and the retention of soil fertility on the plot that is being farmed. Hands had seen the devastating consequences that are caused by slash and burn agriculture while working in Honduras; this new technique seemed to offer the solution to the environmental and economic problems faced by so many slash and burn farmers. Although this technique has the potential to save rainforest and lift many out of poverty, Inga alley cropping has not yet reached its full potential, although the charity Inga Foundation, headed by Mike Hands, has been consulted about potential projects in Haiti ( which is almost completely deforested) and the Congo. Discussions have also been mooted about projects in Peru and Madagascar. Another charity, Rainforest Saver formed to promote Inga Alley Cropping, started a project in 2016 in Ecuador, in the area of the Amazon where Inga edulis originates from, and by the end of 2018 more than 60 farms in the area had Inga plots. Rainforest Saver also started a project in Cameroon in 2009, where in late 2018 there were around 100 farms with Inga plots, mainly in Western Cameroon. Method For Inga alley cropping the trees are planted in rows (hedges) close together, with a gap, the alley, of about 4m between the rows. An initial application of rock phosphate has kept the system going for many years. When the trees have grown, usually in about two years, the canopies close over the alley and cut off the light and so smother the weeds. The trees are then carefully pruned. The larger branches are used for firewood. The smaller branches and leaves are left on the ground in the alleys. These rot down into a good mulch (compost). If any weeds haven't been killed off by lack of light the mulch smothers them. The farmer then pokes holes into the mulch and plants their crops into the holes. The crops grow, fed by the mulch. The crops feed on the lower layers while the latest prunings form a protective layer over the soil and roots, shielding them from both the hot sun and heavy rain. This makes it possible for the roots of both the crops and the trees to stay to a considerable extent in the top layer of soil and the mulch, thus benefiting from the food in the mulch, and escaping soil pests and toxic minerals lower down. Pruning the Inga also makes its roots die back, thus reducing competition with the crops. Forest farming In forest farming, high-value crops are grown under a suitably-managed tree canopy. This is sometimes called multi-story cropping, or in tropical villages as home gardening. It can be practised at varying levels of intensity but always involves some degree of management; this distinguishes it from simple harvesting of wild plants from the forest. Riparian forest buffers Riparian buffers are strips of permanent vegetation located along or near active watercourses or in ditches where water runoff concentrates. The purpose is to keep nutrients and soil from contaminating the water. Silvopasture Trees can benefit fauna in a silvopasture system, where cattle, goats, or sheep browse on grasses grown under trees. In hot climates, the animals are less stressed and put on weight faster when grazing in a cooler, shaded environment. The leaves of trees or shrubs can also serve as fodder. Similar systems support other fauna. Deer and pigs gain when living and feeding in a forest ecosystem, especially when the tree forage nourishes them. In aquaforestry, trees shade fish ponds. In many cases, the fish eat the leaves or fruit from the trees. The dehesa or montado system of silviculture are an example of pigs and bulls being held extensively in Spain and Portugal. Windbreaks Windbreaks reduce wind velocity over and around crops. This increases yields through reduced drying of the crop and/or by preventing the crop from toppling in strong wind gusts. In Switzerland Since the 1950s, four-fifths of Swiss Hochstammobstgärten (traditional orchards with tall trees) have disappeared. An agroforestry scheme was tested here with trees together with annual crops. Trees tested were walnut (Juglans regia) and cherry (Prunus avium). Forty to seventy trees per hectare were recommended, yields were somewhat decreasing with increasing tree height and foliage. However, the total yield per area is shown to be up to 30 percent higher than for monocultural systems. Another set of tests involve growing Populus tremula for biofuel at 52 trees a hectare and with grazing pasture alternated every two to three years with maize or sorghum, wheat, strawberries and fallowing between rows of modern short-pruned & grafted apple cultivars ('Boskoop' & 'Spartan') and growing modern sour cherry cultivars ('Morina', 'Coraline' and 'Achat') and apples, with bushes in the rows with tree (dogrose, Cornus mas, Hippophae rhamnoides) intercropped with various vegetables. Forest gardening Forest gardening is a low-maintenance, sustainable, plant-based food production and agroforestry system based on woodland ecosystems, incorporating fruit and nut trees, shrubs, herbs, vines and perennial vegetables which have yields directly useful to humans. Making use of companion planting, these can be intermixed to grow in a succession of layers to build a woodland habitat. Forest gardening is a prehistoric method of securing food in tropical areas. In the 1980s, Robert Hart coined the term "forest gardening" after adapting the principles and applying them to temperate climates. History Since prehistoric times, hunter-gatherers might have influenced forests, for instance in Europe by Mesolithic people bringing favored plants like hazel with them. Forest gardens are probably the world's oldest form of land use and most resilient agroecosystem. First Nation villages in Alaska with forest gardens filled with nuts, stone fruit, berries, and herbs, were noted by an archeologist from the Smithsonian in the 1930s. Forest gardens are still common in the tropics and known as Kandyan forest gardens in Sri Lanka; , family orchards in Mexico; agroforests; or shrub gardens. They have been shown to be a significant source of income and food security for local populations. Robert Hart adapted forest gardening for the United Kingdom's temperate climate during the 1980s. In temperate climates Hart began farming at Wenlock Edge in Shropshire to provide a healthy and therapeutic environment for himself and his brother Lacon. Starting as relatively conventional smallholders, Hart soon discovered that maintaining large annual vegetable beds, rearing livestock and taking care of an orchard were tasks beyond their strength. However, a small bed of perennial vegetables and herbs he planted was looking after itself with little intervention. Following Hart's adoption of a raw vegan diet for health and personal reasons, he replaced his farm animals with plants. The three main products from a forest garden are fruit, nuts and green leafy vegetables. He created a model forest garden from a 0.12 acre (500 m2) orchard on his farm and intended naming his gardening method ecological horticulture or ecocultivation. Hart later dropped these terms once he became aware that agroforestry and forest gardens were already being used to describe similar systems in other parts of the world. He was inspired by the forest farming methods of Toyohiko Kagawa and James Sholto Douglas, and the productivity of the Keralan home gardens; as Hart explained, "From the agroforestry point of view, perhaps the world's most advanced country is the Indian state of Kerala, which boasts no fewer than three and a half million forest gardens ... As an example of the extraordinary intensity of cultivation of some forest gardens, one plot of only was found by a study group to have twenty-three young coconut palms, twelve cloves, fifty-six bananas, and forty-nine pineapples, with thirty pepper vines trained up its trees. In addition, the smallholder grew fodder for his house-cow." Seven-layer system Further development The Agroforestry Research Trust, managed by Martin Crawford, runs experimental forest gardening projects on a number of plots in Devon, United Kingdom. Crawford describes a forest garden as a low-maintenance way of sustainably producing food and other household products. Ken Fern had the idea that for a successful temperate forest garden a wider range of edible shade tolerant plants would need to be used. To this end, Fern created the organisation Plants for a Future which compiled a plant database suitable for such a system. Fern used the term woodland gardening, rather than forest gardening, in his book Plants for a Future. Kathleen Jannaway, the cofounder of Movement for Compassionate Living (MCL) with her husband Jack, wrote a book outlining a sustainable vegan future called Abundant Living in the Coming Age of the Tree in 1991. The MCL promotes forest gardening and other types of vegan organic gardening. In 2009 it provided a grant of £1,000 to the Bangor Forest Garden project in Gwynedd, North West Wales. Permaculture Bill Mollison, who coined the term permaculture, visited Hart at his forest garden in October 1990. Hart's seven-layer system has since been adopted as a common permaculture design element. Numerous permaculturalists are proponents of forest gardens, or food forests, such as Graham Bell, Patrick Whitefield, Dave Jacke, Eric Toensmeier and Geoff Lawton. Bell started building his forest garden in 1991 and wrote the book The Permaculture Garden in 1995, Whitefield wrote the book How to Make a Forest Garden in 2002, Jacke and Toensmeier co-authored the two volume book set Edible Forest Gardens in 2005, and Lawton presented the film Establishing a Food Forest in 2008. Geographical distribution Forest gardens, or home gardens, are common in the tropics, using intercropping to cultivate trees, crops, and livestock on the same land. In Kerala in south India as well as in northeastern India, the home garden is the most common form of land use and is also found in Indonesia. One example combines coconut, black pepper, cocoa and pineapple. These gardens exemplify polyculture, and conserve much crop genetic diversity and heirloom plants that are not found in monocultures. Forest gardens have been loosely compared to the religious concept of the Garden of Eden. Americas The Amazon rainforest, rather than being a pristine wilderness, has been shaped by humans for at least 11,000 years through practices such as forest gardening and terra preta. Since the 1970s, numerous geoglyphs have been discovered on deforested land in the Amazon rainforest, furthering the evidence of pre-Columbian civilizations. On the Yucatán Peninsula, much of the Maya food supply was grown in "orchard gardens", known as pet kot. The system takes its name from the low wall of stones (pet meaning 'circular' and kot, 'wall of loose stones') that characteristically surrounds the gardens. The environmental historian William Cronon argued in his 1983 book Changes in the Land that indigenous North Americans used controlled burning to form ideal habitat for wild game. The natural environment of New England was sculpted into a mosaic of habitats. When indigenous Americans hunted, they were "harvesting a foodstuff which they had consciously been instrumental in creating". Most English settlers, however, assumed that the wealth of food provided by the forest was a result of natural forces, and that indigenous people lived off "the unplanted bounties of nature." Animal populations declined after settlement, while fields of strawberries and raspberries found by the earliest settlers became overgrown and disappeared for want of maintenance. Plants Some plants, such as wild yam, work as both a root plant and as a vine. Ground covers are low-growing edible forest garden plants that help keep weeds in control and provide a way to utilize areas that would otherwise be unused. Cardamom Ginger Chervil Bergamot Sweet woodruff Sweet cicely Projects El Pilar on the Belize–Guatemala border features a forest garden to demonstrate traditional Maya agricultural practices. A further one acre model forest garden, called Känan K'aax (meaning 'well-tended garden' in Mayan), is funded by the National Geographic Society and developed at Santa Familia Primary School in Cayo. In the United States, the largest known food forest on public land is believed to be the seven acre Beacon Food Forest in Seattle, Washington. Other forest garden projects include those at the central Rocky Mountain Permaculture Institute in Basalt, Colorado, and Montview Neighborhood farm in Northampton, Massachusetts. The Boston Food Forest Coalition promotes local forest gardens. In Canada Richard Walker has been developing and maintaining food forests in British Columbia for over 30 years. He developed a three-acre food forest that at maturity provided raw materials for a plant nursery and herbal business as well as food for his family. The Living Centre has developed various forest garden projects in Ontario. In the United Kingdom, other than those run by the Agroforestry Research Trust (ART), projects include the Bangor Forest Garden in Gwynedd, northwest Wales. Martin Crawford from ART administers the Forest Garden Network, an informal network of people and organisations who are cultivating forest gardens. Since 2014, Gisela Mir and Mark Biffen have been developing a small-scale edible forest garden in Cardedeu near Barcelona, Spain, for experimentation and demonstration. Forest farming Forest farming is the cultivation of high-value specialty crops under a forest canopy that is intentionally modified or maintained to provide shade levels and habitat that favor growth and enhance production levels. Forest farming encompasses a range of cultivated systems from introducing plants into the understory of a timber stand to modifying forest stands to enhance the marketability and sustainable production of existing plants. Forest farming is a type of agroforestry practice characterized by the "four I's": intentional, integrated, intensive and interactive. Agroforestry is a land management system that combines trees with crops or livestock, or both, on the same piece of land. It focuses on increasing benefits to the landowner as well as maintaining forest integrity and environmental health. The practice involves cultivating non-timber forest products or niche crops, some of which, such as ginseng or shiitake mushrooms, can have high market value. Non-timber forest products (NTFPs) are plants, parts of plants, fungi, and other biological materials harvested from within and on the edges of natural, manipulated, or disturbed forests. Examples of crops are ginseng, shiitake mushrooms, decorative ferns, and pine straw. Products typically fit into the following categories: edible, medicinal and dietary supplements, floral or decorative, or specialty wood-based products. History Forest farming, though not always by that name, is practiced around the world. For centuries, humans have relied on fruits, nuts, seeds, parts of foliage and pods from trees and shrubs in the forests to feed themselves and their livestock. Over time, certain species have been selected for cultivation near homes or livestock to provide food or medicine. For example, in the southern United States, mulberry trees are used as a feedstock for pigs and often cultivated near pig quarters. In 1929, J. Russell Smith, Emeritus Professor of Economic Geography at Columbia University, published "Tree Crops – A Permanent Agriculture" which stated that crop-yielding trees could provide useful substitutes for cereals in animal feeding programs, as well as conserve environmental health. Toyohiko Kagawa read and was heavily influenced by Smith’s publication and began experimental cultivation under trees in Japan during the 1930s. Through forest farming, or three-dimensional forestry, Kagawa addressed problems of soil erosion by persuading many of Japan's upland farmers to plant fodder trees to conserve soil, supply food and feed animals. He combined extensive plantings of walnut trees, harvested the nuts and fed them to the pigs, then sold the pigs as a source of income. When the walnut trees matured, they were sold for timber and more trees were planted so that there was a continuous cycle of economic cropping that provided both short-term and long-term income to the small landowner. The success of these trials prompted similar research in other countries. World War II disrupted communication and slowed advances in forest farming. In the mid-1950s research resumed in places such as southern Africa. Kagawa was also an inspiration to Robert Hart pioneered forest gardening in temperate climates in the sixties in Shropshire, England. In earlier years, livestock were often considered part of the forest farming system. Now they are typically excluded and agroforestry systems that integrate trees, forages and livestock are referred to as silvopastures. Because forest farming combines ecological stability of natural forests and productive agriculture systems, it is considered to have great potential for regenerating soils, restoring ground water supplies, controlling floods and droughts and cultivating marginal lands. Principles Forest farming principles constitute an ecological approach to forest management. Forest resources are judiciously used while biodiversity and wildlife habitat are conserved. Forest farms have the potential to restore ecological balance to fragmented second growth forests through intentional manipulation to create the desired forest ecosystem. In some instances, the intentional introduction of species for botanicals, medicinals, food or decorative products is accomplished using existing forests. The tree cover, soil type, water supply, land form and other site characteristics determine what species will thrive. Developing an understanding of species/site relationships as well as understanding the site limitations is necessary to utilize these resources for production needs, while conserving adequate resources for the long-term health of the forest. Apart from the environmental benefits, forest farming can increase the economic value of forest property and provide short- and long-term benefits to the landowner. Forest farming provides economic return from intact forest ecosystems, but timber sales can remain part of the long-term management strategy. Methods Forest farming methods may include: Intensive, yet careful thinning of overstocked, suppressed tree stands; multiple integrated entries to accomplish thinning so that systemic shock is minimized; and interactive management to maintain a cross-section of healthy trees and shrubs of all ages and species. Physical disturbance to the surrounding area should be minimized. The following are forest farming techniques described in the Training Manual produced by the Center for Agroforestry at the University of Missouri. Level of management that is required (from most intense to least intense) 1. Forest gardening is the most intensive of forest farming methods. In addition to thinning the overstory, this method involves clearing the understory of undesirable vegetation and other practices that are closely related to agronomy (tillage, fertilization, weeding, and control of disease and insects and wildlife management). Due to input levels, this method often produces lower valued products compared to other methods. Forest gardens take advantage of the vertical levels of light availability and space under the forest canopy so that more than one crop can be grown at once if desired. 2. Wild-simulated seeks to maintain a natural growing environment, yet enriches local NTFP populations to create an abundant renewable supply of the products. Minimal disturbance and natural growing conditions ensure products will be similar in appearance and quality of those harvested from the wild. Rather than till, practitioners often rake leaves to expose soil, sow seed directly onto the ground, and then cover with leaves again. Since this method produces NTFPs that closely resemble wild plants; they often command a higher price than NTFPs produced using the forest gardening method. 3. Forest tending involves adjusting tree crown density to manipulate light levels that favor natural reproduction of desirable NTFPs. This low intensity management approach does not involve supplemental planting to increase populations of desired NTFPs. 4. Wildcrafting is the harvesting of naturally growing NTFPs. It is not considered a forest farming practice since there is no human involvement in the plant’s establishment and maintenance. However, wildcrafters often take steps to protect NTFPs with future harvests in mind. It becomes agroforestry once forest thinnings, or other inputs, are applied to sustain or maintain plant populations that might otherwise succumb to successional changes in the forest. The most important difference between forest farming and wildcrafting is that forest farming intentionally produces NTFPS, whereas wildcrafting seeks and gathers from naturally growing NTFPs. Production considerations Forest farming can be a small business opportunity for landowners and requires careful planning, including a business and marketing plan. Learning how to market the NTFPs on the Internet is an option, but may entail higher shipping costs. Landowners should consider all options for selling their products including, farmer’s markets or restaurants that focus on locally grown ingredients. The development phase should include a forest management plan that states the landowner’s objectives and a resource inventory. Start-up costs should be analyzed as specific equipment may be necessary to harvest or process the product, whereas other crops require minimal initial investment. Local incentives for sustainable forest management, as well as regulations and policies should be explored. The Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES) regulates international trade of certain plant (American ginseng and goldenseal) and animal species. To be legally exported, regulated plants must be harvested and records kept according to CITES rules and restrictions. Many states also have harvesting regulations for certain native plants that are searchable online. Another good source to start with on information is the Medicinal Plants at Risk 2008 report, by the Center for Biological Diversity] in the U.S. Examples of crops (from the National Agroforestry Center) Medicinal herbs: Ginseng (Panax quinquefolius) Black Cohosh (Actaea racemosa) Goldenseal (Hydrastis canadensis) Bloodroot (Sanguinaria canadensis) Pacific yew (Taxus brevifolia) Mayapple (Podophyllum peltatum) Saw palmetto (Serenoa repens) American Pokeweed (Phytolacca americana) Nuts: Black walnut (Juglans nigra) Hazelnut (Corylus avellana) Shagbark hickory (Carya ovata) Beechnut (Fagus sylvatica) Fruit: Pawpaw (Asimina triloba) Currants (Ribes spp) Elderberry (Sambucus spp) Serviceberry (Amelanchier spp) Blackberry (Rubus spp) Huckleberry (Gaylussacia brachycera) Other food crops: Ramps (wild leeks) (Allium tricoccum) Syrups (maple) Honey Mushrooms Other edible roots Other products: (mulch, decoratives, crafts, dyes) Pine straw Willow twigs Vines Beargrass (Xerophyllum tenax) Ferns Pine cones Moss Native ornamentals: Rhododendron (Rhododendron catawbiense) Highbush cranberry (Viburnum trilobum) Flowering dogwood (Cornus florida) Farmer-managed natural regeneration Farmer-managed natural regeneration (FMNR) is a low-cost, sustainable land restoration technique used to combat poverty and hunger amongst poor subsistence farmers in developing countries by increasing food and timber production, and resilience to climate extremes. It involves the systematic regeneration and management of trees and shrubs from tree stumps, roots and seeds. FMNR was developed by the Australian agricultural economist Tony Rinaudo in the 1980s in West Africa. The background and development are described in Rinaudo's book The Forest Underground. FMNR is especially applicable, but not restricted to, the dryland tropics. As well as returning degraded croplands and grazing lands to productivity, it can be used to restore degraded forests, thereby reversing biodiversity loss and reducing vulnerability to climate change. FMNR can also play an important role in maintaining not-yet-degraded landscapes in a productive state, especially when combined with other sustainable land management practices such as conservation agriculture on cropland and holistic management on range lands. FMNR adapts centuries-old methods of woodland management, called coppicing and pollarding, to produce continuous tree-growth for fuel, building materials, food and fodder without the need for frequent and costly replanting. On farmland, selected trees are trimmed and pruned to maximise growth while promoting optimal growing conditions for annual crops (such as access to water and sunlight). When FMNR trees are integrated into crops and grazing pastures there is an increase in crop yields, soil fertility and organic matter, soil moisture and leaf fodder. There is also a decrease in wind and heat damage, and soil erosion. FMNR complements the evergreen agriculture, conservation agriculture and agroforestry movements. It is considered a good entry point for resource-poor and risk-averse farmers to adopt a low-cost and low-risk technique. This in turn has acted as a stepping stone to greater agricultural intensification as farmers become more receptive to new ideas. Background Throughout the developing world, immense tracts of farmland, grazing lands and forests have become degraded to the point they are no longer productive. Deforestation continues at a rapid pace. In Africa's drier regions, 74 percent of rangelands and 61 percent of rain-fed croplands are damaged by moderate to very severe desertification. In some African countries deforestation rates exceed planting rates by 300:1. Degraded land has an extremely detrimental effect on the lives of subsistence farmers who depend on it for their food and livelihoods. Subsistence farmers often make up to 70–80 percent of the population in these regions and they regularly suffer from hunger, malnutrition and even famine as a consequence. In the Sahel region of Africa, a band of savanna which runs across the continent immediately south of the Sahara Desert, large tracts of once-productive farmland are turning to desert. In tropical regions across the world, where rich soils and good rainfall would normally assure bountiful harvests and fat livestock, some environments have become so degraded they are no longer productive. Severe famines across the African Sahel in the 1970s and 1980s led to a global response, and stopping desertification became a top priority. Conventional methods of raising exotic and indigenous tree species in nurseries were used. Despite investing millions of dollars and thousands of hours of labour, there was little overall impact. Conventional approaches to reforestation in such harsh environments faced insurmountable problems and were costly and labour-intensive. Once planted out, drought, sand storms, pests, competition from weeds and destruction by people and animals negated efforts. Low levels of community ownership were another inhibiting factor. Existing indigenous vegetation was generally dismissed as 'useless bush', and it was often cleared to make way for exotic species. Exotics were planted in fields containing living and sprouting stumps of indigenous vegetation, the presence of which was barely acknowledged, let alone seen as important. This was an enormous oversight. In fact, these living tree stumps are so numerous they constitute a vast 'underground forest' just waiting for some care to grow and provide multiple benefits at little or no cost. Each stump can produce between 10 and 30 stems each. During the process of traditional land preparation, farmers saw the stems as weeds and slashed and burnt them before sowing their food crops. The net result was a barren landscape for much of the year with few mature trees remaining. To the casual observer, the land was turning to desert. Most concluded that there were no trees present and that the only way to reverse the problem was through tree planting. Meanwhile, established indigenous trees continued to disappear at an alarming rate. In Niger, from the 1930s until 1993, forestry laws took tree ownership and responsibility for the care of trees out of the hands of the people. Reforestation through conventional tree planting seemed to be the only way to address desertification at the time. History In the early-1980s, in the Maradi region of the Republic of Niger, the missionary organisation, Serving in Mission (SIM), was unsuccessfully attempting to reforest the surrounding districts using conventional means. In 1983, SIM began experimenting and promoting FMNR amongst about 10 farmers. During the famine of 1984, a food-for-work program was introduced that saw some 70,000 people exposed to FMNR and its practice on around 12,500 hectares of farmland. From 1985 to 1999, FMNR continued to be promoted locally and nationally as exchange visits and training days were organised for various NGOs, government foresters, Peace Corps volunteers, and farmer and civil society groups. Additionally, SIM project staff and farmers visited numerous locations across Niger to provide training. By 2004 it was ascertained that FMNR was being practised on over five million hectares or 50 percent of Niger's farmland – an average reforestation rate of 250,000 hectares per year over a 20-year period. This transformation prompted a Senior Fellow of the World Resources Institute, Chris Reij, to comment that "this is probably the largest positive environmental transformation in the Sahel and perhaps all of Africa". In 2004, World Vision Australia and World Vision Ethiopia initiated a forestry-based carbon sequestration project as a potential means to stimulate community development while engaging in environmental restoration. A partnership with the World Bank, the Humbo Community-based Natural Regeneration Project involved the regeneration of 2,728 hectares of degraded native forests. This brought social, economic and ecological benefits to the participating communities. Within two years, communities were collecting wild fruits, firewood, and fodder, and reported that wildlife had begun to return and erosion and flooding had been reduced. In addition, the communities are now receiving payments for the sale of carbon credits through the Clean Development Mechanism (CDM) of the Kyoto Protocol. Following the success of the Humbo project, FMNR spread to the Tigray region of northern Ethiopia where 20,000 hectares have been set aside for regeneration, including 10 hectare FMNR model sites for research and demonstration in each of 34 sub-districts. The Government of Ethiopia has committed to reforest 15 million hectares of degraded land using FMNR as part of a climate change and renewable energy plan to become carbon neutral by 2025. In Talensi, northern Ghana, FMNR is being practiced on 2,000–3,000 hectares and new projects are introducing FMNR into three new districts. In the Kaffrine and Diourbel regions of Senegal, FMNR has spread across 50,000 hectares in four years. World Vision is also promoting FMNR in Indonesia, Myanmar and East Timor. There are also examples of both independently promoted and spontaneous FMNR movements occurring. In Burkina Faso, for example, an increasing part of the country is being transformed into agro-forestry parkland. And in Mali, an ageing agro-forestry parkland of about six million hectares is showing signs of regeneration. Key principles FMNR depends on the existence of living tree stumps or roots in crop fields, grazing pastures, woodlands or forests. Each season bushy growth will sprout from the stumps/roots often appearing like small shrubs. Continuous grazing by livestock, regular burning and/or regular harvesting for fuel wood results in these 'shrubs' never attaining tree stature. On farmland, standard practice has been for farmers to slash this regrowth in preparation for planting crops, but with a little attention this growth can be turned into a valuable resource without jeopardising crop yields. For each stump, a decision is made as to how many stems will be chosen to grow. The tallest and straightest stems are selected and the remaining stems culled. Best results are obtained when the farmer returns regularly to prune any unwanted new stems and side branches as they appear. Farmers can then grow other crops between and around the trees. When farmers want wood they can cut the stem(s) they want and leave the rest to continue growing. The remaining stems will increase in size and value each year, and will continue to protect the environment. Each time a stem is harvested, a younger stem is selected to replace it. Various naturally occurring tree species can be used which may also provide berries, fruits and nuts or have medicinal qualities. In Niger, commonly used species include: Strychnos spinosa, Balanites aegyptiaca, Boscia senegalensis, Ziziphus spp., Annona senegalensis, Poupartia birrea and Faidherbia albida. However, the most important determinants are whatever species are locally available, their ability to re-sprout after cutting, and the value local people place on those species. Faidherbia albida, also known as the 'fertiliser tree', is popular for intercropping across the Sahel as it fixes nitrogen into the soil, provides fodder for livestock, and shade for crops and livestock. By shedding its leaves in the wet season, Faidherbia provides beneficial light shade to crops when high temperatures would otherwise damage crops or retard growth. Leaf fall contributes useful nutrients and organic matter to the soil. The practice of FMNR is not confined to croplands. It is being practised on grazing land and in degraded communal forests as well. When there are no living stumps, seeds of naturally occurring species are used. In reality, there is no fixed way of practising FMNR and farmers are free to choose which species they will leave, the density of trees they prefer, and the timing and method of pruning. In practice Benefits FMNR can restore degraded farmlands, pastures and forests by increasing the quantity and value of woody vegetation, by increasing biodiversity and by improving soil structure and fertility through leaf litter and nutrient cycling. The reforestation also retards wind and water erosion; it creates windbreaks which decrease soil moisture evaporation, and protects crops and livestock against searing winds and temperatures. Often, dried up springs reappear and the water table rises towards historic levels; insect eating predators including insects, spiders and birds return, helping to keep crop pests in check; the trees can be a source of edible berries and nuts; and over time the biodiversity of plant and animal life is increased. FMNR can be used to combat deforestation and desertification and can also be an important tool in maintaining the integrity and productivity of land that is not yet degraded. Trials, long-running programs and anecdotal data indicate that FMNR can at least double crop yields on low fertility soils. In the Sahel, high numbers of livestock and an eight month dry season can mean that pastures are completely depleted before the rains commence. However, with the presence of trees, grazing animals can make it through the dry season by feeding on tree leaves and seed pods of some species, at a time when no other fodder is available. In northeast Ghana, more grass became available with the introduction of FMNR because communities worked together to prevent bush fires from destroying their trees. Well designed and executed FMNR projects can act as catalysts to empower communities as they negotiate land ownership or user rights for the trees in their care. This assists with self-organisation, and with the development of new agriculture-based micro-enterprises (e.g., selling firewood, timber and handcrafts made from timber or woven grasses). Conventional approaches to reversing desertification, such as funding tree planting, rarely spread beyond the project boundary once external funding is withdrawn. By comparison, FMNR is cheap, rapid, locally led and implemented. It uses local skills and resources – the poorest farmers can learn by observation and teach their neighbours. Given an enabling environment, or at least the absence of a 'disabling' environment, FMNR can be done at scale and spread well beyond the original target area without ongoing government or NGO intervention. World Vision evaluations of FMNR conducted in Senegal and Ghana in 2011 and 2012 found that households practising FMNR were less vulnerable to extreme weather shocks such as drought and damaging rain and wind storms. The following table summarises FMNR's benefits which fit the sustainable development model of economic, social and environmental benefits: Sources: Key success factors and constraints While there are numerous accounts of the uptake and spread of FMNR independent of aid and development agencies, the following factors have been found to be beneficial for its introduction and spread: Awareness creation of FMNR's potential. Capacity building through workshops and exchange visits. Awareness of the devastating effects of deforestation. The adoption of FMNR is more likely when communities acknowledge their situation and the need to take action. This perception of need can be supported by education. An FMNR champion/facilitator from within the community who encourages, challenges and trains peers. This is critical during the first three to five years, and continues to be important for up to 10 years. Regular site visits also ensure early detection and remedial action on resistance and threats to FMNR through deliberate damage to trees and theft. The buy-in of all stakeholders including their agreement on any by-laws created for FMNR and the consequences for infringements. Stakeholders include FMNR practitioners, local, regional and national government departments of agriculture and forestry, men, women, youth, marginalised groups (including nomadic herders), cultivators and commercial interests. Stakeholder buy-in is also important to create a critical mass of FMNR adopters in order to change social attitudes from a position of apathy or active participation in deforestation to one of proactive sustainable tree management through FMNR. Government support through the creation of favourable policies, positive reinforcement of actions facilitating the spread of FMNR, and disincentives for actions working against the spread of FMNR. FMNR practitioners need to be confident that they will benefit from their labours (either private or community ownership of trees, or legally binding user rights). Reinforcement of existing organisational structures (farmers clubs, development groups, traditional leadership structures) or establishment of new structures which will provide a framework for communities to practise FMNR on a local, district or region-wide basis. A communications strategy which includes education in schools, radio programs and engagement with religious and traditional leaders to become advocates. Establishment of a legal, transparent and accessible market for FMNR wood and non-timber forest products, enabling practitioners to benefit financially from their activities. Brown et al. suggest that the two main reasons why FMNR has spread so widely in Niger are attitudinal change by the community of what constitutes good land management practices, and farmers' ownership of trees. Farmers need the assurance that they will benefit from their labour. Giving farmers either outright ownership of the trees they protect, or tree-user rights, has made it possible for large-scale farmer-led reforestation to take place. Current and future directions Over nearly 30 years, FMNR has changed the farming landscape in some of the poorest countries in the world, including parts of Niger, Burkina Faso, Mali, and Senegal, providing subsistence farmers with the methods necessary to become more food secure and resilient against severe weather events. The 2011–2012 food crisis in East Africa gave a stark reminder of the importance of addressing root causes of hunger. In the 2011 State of the World Report, Bunch concludes that four major factors – lack of sustainable fertile land, loss of traditional fallowing, cost of fertiliser and climate change – are coming together all at once in a sort of "perfect storm" that will almost surely result in an African famine of unprecedented proportions, probably within the next four to five years. It will most heavily affect the lowland, semi-arid to sub-humid areas of Africa (including the Sahel, parts of eastern Africa, plus a band from Malawi across to Angola and Namibia); and unless the world does something dramatic, 10 to 30 million people could die from famine between 2015 and 2020. Restoration of degraded land through FMNR is one way of addressing these major contributors to hunger. In recent years FMNR has come to the attention of global development agencies and grassroots movements alike. The World Bank, World Resources Institute, World Agroforestry Center, USAID and the Permaculture movement are amongst those either actively promoting or advocating for the uptake of FMNR and FMNR has received recognition from a number of quarters including: In 2010, FMNR won the Interaction 4 Best Practice and Innovation Initiative award in recognition of high technical standards and effectiveness in addressing the food security and livelihood needs of small producers in the areas of natural resource management and agro forestry. In 2011, FMNR won the World Vision International Global Resilience Award for the most innovative initiative in the area of resilient development practice and natural environment and climate issues. In 2012 WVA was awarded the Arbor Day Award for Education Innovation. In April 2012, World Vision Australia – in partnership with the World Agroforestry Center and World Vision East Africa – held an international conference in Nairobi called "Beating Famine" to analyse and plan how to improve food security for the world's poor through the use of FMNR and Evergreen Agriculture. The conference was attended by more than 200 participants, including world leaders in sustainable agriculture, five East African ministers of agriculture and the environment, ambassadors, and other government representatives from Africa, Europe, and Australia, and leaders from non-government and international organisations. Two major outcomes of the conference were: The establishment of a global FMNR network of key stakeholders to promote, encourage and initiate the scale-up of FMNR globally. Country, regional and global level plans as a basis for inter-organisation collaboration for FMNR scale-up. The conference acted as a catalyst for media coverage of FMNR in some of the world's leading outlets and a noticeable increase in momentum for an FMNR global movement. This heightened awareness of FMNR has created an opportunity for it to spread exponentially worldwide.	Technology	Trees and forestry	null
1489133	https://en.wikipedia.org/wiki/Nephrops%20norvegicus	Nephrops norvegicus	Nephrops norvegicus, known variously as the Norway lobster, Dublin Bay prawn, (compare langostino) or scampi, is a slim, coral colored lobster that grows up to long, and is "the most important commercial crustacean in Europe". It is now the only extant species in the genus Nephrops, after several other species were moved to the closely related genus Metanephrops. It lives in the north-eastern Atlantic Ocean, and parts of the Mediterranean Sea, but is absent from the Baltic Sea and Black Sea. Adults emerge from their burrows at night to feed on worms and fish. Description Nephrops norvegicus has the typical body shape of a lobster, albeit narrower than the large genus Homarus. It is pale orange in colour, and grows to a typical length of , or exceptionally long, including the tail and claws. A carapace covers the animal's cephalothorax, while the abdomen is long and segmented, ending in a broad tail fan. The first three pairs of legs bear claws, of which the first are greatly elongated and bear ridges of spines. Of the two pairs of antennae, the second is the longer and thinner. There is a long, spinous rostrum, and the compound eyes are kidney-shaped, providing the name of the genus, from the Greek roots (nephros, "kidney") and ὄψ ("eye"). Distribution Nephrops norvegicus is found in the north-eastern Atlantic Ocean and North Sea as far north as Iceland and northern Norway, and south to Portugal. It is found in the Mediterranean Sea and is common in the Adriatic Sea, notably the north Adriatic. It is absent from both the Black Sea and the Baltic Sea. Due to its ecological demands for particular sediments, N. norvegicus has a very patchy distribution, and is divided into over 30 populations. These populations are separated by inhospitable terrain, and adults rarely travel distances greater than a few hundred metres. Ecology Nephrops norvegicus adults prefer to inhabit muddy seabed sediments, with more than 40 percent silt and clay. Their burrows are semi-permanent, and vary in structure and size. Typical burrows are deep, with a distance of between the front and back entrances. Norway lobsters spend most of their time either lying in their burrows or by the entrance, leaving their shelters only to forage or mate. Diet Nephrops norvegicus is a scavenger and predator that makes short foraging excursions, mainly during periods of subdued light. They feed on active prey, including worms and fish, which they capture with their chelipeds and walking legs, and food is conveyed to the mouth using the anterior walking legs, assisted by the maxillipeds. There is evidence that Nephrops norvegicus is a major eater of jellyfish. Parasites and symbionts Nephrops norvegicus is the host to a number of parasites and symbionts. A number of sessile organisms attach to the exoskeleton of N. norvegicus, including the barnacle Balanus crenatus and the foraminiferan Cyclogyra, but overall Nephrops suffers fewer infestations of such epibionts than other decapod crustaceans do. In December 1995, the commensal Symbion pandora was discovered attached to the mouthparts of Nephrops norvegicus, and was found to be the first member of a new phylum, Cycliophora, a finding described by Simon Conway Morris as "the zoological highlight of the decade". S. pandora has been found in many populations of N. norvegicus, both in the north Atlantic and in the Mediterranean Sea. Individuals may be found on most segments of the lobster's mouthparts, but are generally concentrated on the central parts of the larger mouthparts, from the mandible to the third maxilliped. The most significant parasite of N. norvegicus is a dinoflagellate of the genus Hematodinium, which has caused epidemic infection in fished populations of N. norvegicus since the 1980s. Hematodinium is a genus that contains major pathogens of a wide variety of decapod crustaceans, although its internal taxonomy is poorly resolved. The species which attacks N. norvegicus causes a syndrome originally described as "post-moult syndrome", in which the carapace turns opaque and becomes highly pigmented, the haemolymph becomes milky white, and the animal appears moribund. Other parasites of N. norvegicus include the gregarine protozoan Porospora nephropis, the trematode Stichocotyle nephropis and the polychaete Histriobdella homari. Life cycle The typical life span of N. norvegicus is 5–10 years, reaching 15 years in exceptional cases. Its reproductive cycle varies depending on geographical position: "the periods of hatching and spawning, and the length of the incubation period, vary with latitude and the breeding cycle changes from annual to biennial as one moves from south to north". Incubation of eggs is temperature-dependent, and in colder climates, the duration of the incubation period increases. This means that, by the time hatching occurs, it may be too late for the females to take part in that year's breeding cycle. In warmer climates, the combined effects of recovery from moulting and ovary maturation mean that spawning can become delayed. This, in turn, has the effect of the female missing out a year of egg carrying. Adult male Nephrops norvegicus moult once or twice a year (usually in late winter or spring) and adult females moult up to once a year (in late winter or spring, after hatching of the eggs). In annual breeding cycles, mating takes place in the spring or winter, when the females are in the soft, post-moult state. The ovaries mature throughout the spring and summer months, and egg-laying takes place in late summer or early autumn. After spawning, the berried (egg-carrying) females return to their burrows and remain there until the end of the incubation period. Hatching takes place in late winter or early spring. Soon after hatching, the females moult and mate again. During the planktonic larval stage (typically 1 to 2 months in duration) the nephrops larvae exhibit a diel vertical migration behaviour as they are dispersed by the local currents. This complex biophysical interaction determines the fate of the larvae; the overlap between advective pathway destination and spatial distributions of suitable benthic habitats must be favourable in order for the larvae to settle and reach maturity. Fisheries The muscular tail of Nephrops norvegicus is frequently eaten, and its meat is known as scampi or langoustine. N. norvegicus is eaten only on special occasions in Spain and Portugal, where it is less expensive than the common lobster, Homarus gammarus. N. norvegicus is an important species for fisheries, being caught mostly by trawling. Around 60,000 tonnes are caught annually, half of it in the United Kingdom's waters. Besides the established trawling fleets, a significant number of fleets using lobster creels have developed. The better size and condition of lobsters caught by this method yield prices three to four times higher than animals netted by trawling. Creel fishing was found to have a reduced impact on the seafloor, require lower fuel consumption, and allow fishermen with smaller boats to participate in this high-value fishery. It has therefore been described as a reasonable alternative to demersal towed gears, and the allocation of additional fishing rights for this type of take has been suggested. The North East Atlantic individual biological stocks of Nephrops are identified as functional units. A number of functional units make up the sea areas over which a total allowable catch (TAC) is set annually by the EU Council of Ministers. For example, the TAC set for North Sea Nephrops is based on the aggregate total tonnage of removals recommended by science for nine separate functional unit areas. This method has attracted criticism because it can promote the overexploitation of a specific functional unit even though the overall TAC is under-fished. In 2016, the UK implemented a package of emergency technical measures with the cooperation of the fishing industry aimed at reducing fishing activity to induce recovery of the Nephrops stock in the Farn(e) Deeps off North East England which was close to collapse. A stock assessment completed in 2018 by the International Council for the Exploration of the Sea (ICES) shows that fishing pressure has been cut and this stock is now below FMSY and that stock size is above MSY Btrigger meaning that the Farne Deeps nephrops stock is being fished at a sustainable level. However, ICES also warn that any substantial transfer of the current surplus fishing opportunities from other functional units to the Farne Deeps would rapidly lead to overexploitation. This suggests that controls on fishing effort should continue at least until the biomass reaches a size that is sustainable when measured against the level of fishing activity by all fishermen wanting to target the stock. In July 2023 the area north-east of Farnes Deep was one of three sites designated as a Highly Protected Marine Area. Discards from Nephrops fishery may account for up to 37% of the energy requirements of certain marine scavengers, such as the hagfish Myxine glutinosa. Boats involved in Nephrops fishery also catch a number of fish species such as plaice and sole, and it is thought that without that revenue, Nephrops fishery would be economically unviable. Taxonomic history Nephrops norvegicus was one of the species included by Carl Linnaeus in his 1758 10th edition of , the starting point for zoological nomenclature. In that work, it was listed as , with a type locality of ("in the Norwegian sea"). In choosing a lectotype, Lipke Holthuis restricted the type locality to the Kattegat at the Kullen Peninsula in southern Sweden (). Two synonyms of the species have been published – "Astacus rugosus", described by the eccentric zoologist Constantine Samuel Rafinesque in 1814 from material collected in the Mediterranean Sea, and "Nephropsis cornubiensis", described by Charles Spence Bate and Joshua Brooking Rowe in 1880. As new genera were erected, the species was moved, reaching its current position in 1814, when William Elford Leach erected the genus Nephrops to hold this species alone. Seven fossil species have since been described in the genus. Populations in the Mediterranean Sea are sometimes separated as "Nephrops norvegicus var. meridionalis Zariquiey, 1935", although this taxon is not universally considered valid.	Biology and health sciences	Crayfishes and lobsters	Animals
1489218	https://en.wikipedia.org/wiki/Porcelain%20crab	Porcelain crab	Porcelain crabs are decapod crustaceans in the widespread family Porcellanidae, which superficially resemble true crabs. They have flattened bodies as an adaptation for living in rock crevices. They are delicate, readily losing limbs when attacked, and use their large claws for maintaining territories. They first appeared in the Tithonian age of the Late Jurassic epoch, 145–152 million years ago. Description Porcelain crabs are small, usually with body widths less than . They share the general body plan of a squat lobster, but their bodies are more compact and flattened, an adaptation for living and hiding under rocks. Porcelain crabs are quite fragile animals, and often shed their limbs to escape predators, hence their name. The lost appendage can grow back over several moults. Porcelain crabs have large chelae (claws), which are used for territorial struggles, but not for catching food. The fifth pair of pereiopods is reduced and used for cleaning. Evolution Porcelain crabs are an example of carcinisation, whereby a noncrab-like animal (in this case a relative of a squat lobster) evolves into an animal that resembles a true crab. Porcelain crabs can be distinguished from true crabs by the apparent number of walking legs (three instead of four pairs; the fourth pair is reduced and held against the carapace), and the long antennae originating on the front outside of the eyestalks. The abdomen of the porcelain crab is long and folded underneath it, free to move. Biogeography and ecology Porcelain crabs live in all the world's oceans, except the Arctic Ocean and the Antarctic. They are common under rocks, and can often be found and observed on rocky beaches and shorelines, startled creatures scurrying away when a stone is lifted. They feed by combing plankton and other organic particles from the water using long setae (feathery hair- or bristle-like structures) on the mouthparts. Some of the common species of porcelain crabs in the Caribbean Sea are Petrolisthes quadratus, found in large numbers under rocks in the intertidal, and the red-and-white polka-dotted Porcellana sayana, which lives commensally within the shells inhabited by large hermit crabs. In Hong Kong, Petrolisthes japonicus is common. Diversity , some 4723 extant species of porcelain crab had been described, divided among these 30 genera: Aliaporcellana Nakasone & Miyake, 1969 Allopetrolisthes Haig, 1960 Ancylocheles Haig, 1978 Capilliporcellana Haig, 1978 Clastotoechus Haig, 1960 Enosteoides Johnson, 1970 Euceramus Stimpson, 1860 Eulenaios Ng & Nakasone, 1993 Heteropolyonyx Osawa, 2001 Heteroporcellana Haig, 1978 Liopetrolisthes Haig, 1960 Lissoporcellana Haig, 1978 Madarateuchus Harvey, 1999 Megalobrachium Stimpson, 1858 Minyocerus Stimpson, 1858 Neopetrolisthes Miyake, 1937 Neopisosoma Haig, 1960 Novorostrum Osawa, 1998 Orthochela Glassell, 1936 Pachycheles Stimpson, 1858 Parapetrolisthes Haig, 1962 Petrocheles Miers, 1876 Petrolisthes Stimpson, 1858 Pisidia Leach, 1820 Polyonyx Stimpson, 1858 Porcellana Lamarck, 1801 Porcellanella White, 1852 Pseudoporcellanella Sankarankutty, 1962 Raphidopus Stimpson, 1858 Ulloaia Glassell, 1938 The fossil record of porcelain crabs includes species of Pachycheles, Pisidia, Polyonyx, Porcellana, and a further six genera known only from fossils: Annieporcellana Fraaije et al., 2008 Beripetrolisthes De Angeli & Garassino, 2002 Eopetrolisthes De Angeli & Garassino, 2002 Lobipetrolisthes De Angeli & Garassino, 2002 Longoporcellana Müller & Collins, 1991 The earliest claimed porcelain crab fossil was Jurellana from the Tithonian aged Ernstbrunn Limestone of Austria. However, it was subsequently determined to be a true crab. With the new oldest porcelain crab being Vibrissalana from the same locality.	Biology and health sciences	Crabs and hermit crabs	Animals
2134303	https://en.wikipedia.org/wiki/Birdwing	Birdwing	Birdwings are butterflies in the swallowtail family, that belong to the genera Trogonoptera, Troides, and Ornithoptera. Most recent authorities recognise 36 species, however, this is debated, and some authorities include additional genera. Birdwings are named for their exceptional size, angular wings, and birdlike flight. They are found across tropical Asia, mainland and archipelagic Southeast Asia, and Australasia. Included among the birdwings are some of the largest butterflies in the world: the largest, Queen Alexandra's birdwing; the second largest, the Goliath birdwing; the largest butterfly endemic to Australia, the Cairns birdwing; and the largest butterfly in India, the southern birdwing. Another well-known species is Rajah Brooke's birdwing, a particularly attractive species named after Sir James Brooke, the first White Rajah of 19th-century Sarawak. Due to their size and brightly coloured males, they are popular among collectors of butterflies, but all birdwings are now listed by CITES, thereby limiting (and in the case of O. alexandrae completely banning) international trade. Taxonomy Genera and species genus: Troides subgenus: Ripponia Troides hypolitus – Rippon's birdwing subgenus: Troides species group: Troides aeacus Troides aeacus – golden birdwing Troides dohertyi – Talaud black birdwing Troides magellanus – Magellan birdwing Troides minos – southern birdwing Troides plateni – Dr. Platen's birdwing Troides prattorum – Buru opalescent birdwing Troides rhadamantus – golden birdwing species group: Troides amphrysus Troides amphrysus – Malay birdwing Troides andromache – Borneo birdwing Troides cuneifera Troides miranda – Miranda birdwing species group: Troides haliphron Troides criton – Criton birdwing Troides darsius – Sri Lankan birdwing Troides haliphron – haliphron birdwing Troides plato – silver birdwing Troides riedeli – Riedel's birdwing Troides staudingeri Troides vandepolli – van de Poll's birdwing species group: Troides helena Troides helena – common birdwing Troides oblongomaculatus – oblong-spotted birdwing genus: Trogonoptera Trogonoptera brookiana – Rajah Brooke's birdwing Trogonoptera trojana – Palawan birdwing genus: Ornithoptera subgenus: Aetheoptera Ornithoptera victoriae – Queen Victoria's birdwing subgenus: Ornithoptera Ornithoptera aesacus – Obi Island birdwing Ornithoptera croesus – Wallace's golden birdwing Ornithoptera euphorion – Cairns birdwing Ornithoptera priamus – common green birdwing Ornithoptera richmondia – Richmond birdwing subgenus: Schoenbergia Ornithoptera chimaera – chimaera birdwing Ornithoptera goliath – Goliath birdwing Ornithoptera meridionalis – southern tailed birdwing Ornithoptera paradisea – paradise birdwing Ornithoptera rothschildi – Rothschild's birdwing Ornithoptera tithonus – Tithonus birdwing subgenus: Straatmana Ornithoptera alexandrae – Queen Alexandra's birdwing Natural hybrids Troides prattorum × Troides oblongomaculatus bouruensis — Troides mixtum Ornithoptera rothschildi × Ornithoptera priamus poseidon — Ornithoptera akakeae Ornithoptera victoriae × Ornithoptera priamus urvillianus — Ornithoptera allotei Description Ova After mating, females immediately begin to seek appropriate host plants; climbing vines of the genera Aristolochia and Pararistolochia (both in the family Aristolochiaceae) are sought exclusively. The female lays her spherical eggs under the tips of the vine's leaves, one egg per leaf. Larva The caterpillars are voracious eaters but move very little; a small group will defoliate an entire vine. If starved due to overcrowding, the caterpillars may resort to cannibalism. Fleshy spine-like tubercles line the caterpillars' backs, and their bodies are dark red to brown and velvety black. Some species have tubercles of contrasting colours, often red, or pale "saddle" markings. Like other members of their family, birdwing caterpillars possess a retractable organ behind their heads called an osmeterium. Shaped like the forked tongue of a snake, the osmeterium excretes a fetid terpene-based compound and is deployed when the caterpillar is provoked. The caterpillars are also unappealing to most predators due to their toxicity: the vines which the caterpillars feed upon contain aristolochic acid, a poisonous compound known to be carcinogenic in rats. The feeding caterpillars incorporate and concentrate the aristolochic acid into their tissues, where the poison will persist through metamorphosis and into adulthood. Pupa Birdwing chrysalids are camouflaged to look like a dead leaf or twig. Before pupating, the caterpillars may wander considerable distances from their host plants. In O. alexandrae, it takes about four months to get from egg to adult. Barring predation, this species can also survive up to three months as an adult. Imago Birdwings inhabit rainforests and adults are usually glimpsed along the forest periphery. They feed upon—and are important long-range pollinators of—nectar-bearing flowers of the forest canopy, as well as terrestrial flowers, such as lantana. They are strong flyers and seek sunlit spots in which to bask. Breeding behaviour varies little between species; the female's role is relatively passive, slowly fluttering from perch to perch while the male performs an elaborate, quivering yet stationary dance 20–50 cm above her. Birdwings are typified by large size (up to a maximum body length of 7.6 cm or 3 inches and a wingspan of 28 cm or 11 inches in O. alexandrae), showy colouration (in contrasting shades of green, yellow, black, white, and sometimes blue or orange), and slender, lanceolate forewings. With few exceptions (i.e., the New Guinean O. meridionalis and O. paradisea), the hindwings lack tails. Sexual dimorphism is strong in Ornithoptera species only, where males are black combined with bright iridescent green, blue, orange, or yellow while the larger and less colourful females are overall black or dark brownish with white, pale brown, or yellow markings. Males and females of most Troides birdwings are similar and have jet black to brown dorsal forewings, often with the veins bordered in grey to creamy white. At least one of these darkly-coloured species (T. rhadamantus) possesses thermoreceptors on the anal veins (A2 and A3) of the wings and on the antennal clubs. The antennal receptors of the clubs—which also possess hygroreceptors that measure atmospheric humidity—are known as sensilla basiconica. The thermoreceptors are sensitive to sudden increases in temperature; they are thought to help the butterfly thermoregulate and avoid overheating while basking. The colours of most species are pigmentary (via papiliochrome); but two species, Troides magellanus and the much rarer T. prattorum, are noted for their use of limited-view iridescence: the yellow of the dorsal hindwings is modified by bright blue-green iridescence which is only seen when the butterfly is viewed at a narrow, oblique angle. This "grazing iridescence" is brought about through diffraction of light (after back-reflection) by the wings' extremely steeply-set, multilayered rib-like scales (rather than the ridge-lamellae of most other iridescent butterflies, such as Morpho species). Such limited-view iridescence was previously only known from one other species, the riodinid Ancyluris meliboeus. In A. meliboeus, however, the iridescence is produced by ridge-lamellar scales and features a wider range of colours. The close evolutionary relationship between Troides and Ornithoptera butterflies is well demonstrated by the fact that commercial breeders have produced numerous hybrids between the two. The final and smallest genus is Trogonoptera with just two species. They resemble each other, being overall black with iridescent green markings and a red head. Females are duller than males. Distribution Birdwings are generally found from Southeast Asia to northern Australasia. Trogonoptera brookiana inhabits the Thai-Malay Peninsula, Borneo, Natuna, Sumatra, and various surrounding islands. Trogonoptera trojana is endemic to Palawan in the Philippines. Troides species are distributed widely across the Indomalayan realm, but may be found as far east as New Guinea in the case of Troides oblongomaculatus. Some species may be found as far west as India, and are the westernmost distributed of all birdwings. All Ornithoptera species are found in the northern portion of the Australasian realm, east of Weber's line; the Moluccas, New Guinea, the Solomon Islands, and northeastern Australia. An outlier is Ornithoptera richmondia, which may be found in far northeastern New South Wales, Australia in the southernmost area of its range; the southernmost distribution of all birdwings. Status and protection With the exception of Queen Alexandra's birdwing (O. alexandrae), all birdwings are listed in Appendix II of CITES, and accordingly their trade is restricted in countries that have signed the CITES convention. Exceptions are made for captive-reared specimens, which mainly originate from ranches in Papua New Guinea and Indonesia. Most species of all three genera have now been reared in captivity, though with significant differences in the quantities reared of each species. O. alexandrae is listed on Appendix I and therefore cannot legally be traded internationally. At the 2006 meeting of the CITES Animals Committee some suggested O. alexandrae should be moved to Appendix II, as the conservation benefits of sustainable management perhaps are higher than those of the trade ban. Three Troides and eight Ornithoptera species have been given assessments by the IUCN Red List, with classifications ranging from "least concern" to "endangered". Richmond birdwings (O. richmondia) depend on the plant Aristolochia praevenosa which they need for their caterpillars. However, the very similar Aristolochia elegans (Dutchman's pipe) which can be found in many Australian backyards, kills the caterpillars. Reproduction Ornithoptera, or the genus of birdwing butterflies, usually reproduce sexually and are oviparous. In butterflies sex is determined by a WW/WZ system, with a heterogametic female, reverse of that found in mammals and many other insects, which have a heterogametic male. During copulation males will transfer an ejaculate containing both sperm and accessory substances that can make up to fifteen percent of a males body mass. Mating systems Mating systems, first explored in evolutionary terms by Darwin, includes all behaviours associated with sexual reproduction. Mating systems include all costs and benefits, pre- and postcopulatory competitions, displays and mate choice. Butterfly mating systems have great variation, including strict monandry, one male and one female, to polyandry, having many mates of the opposite sex. Typically Ornithoptera tend to be polygamous, mating with more than one individual. Female choice Female choice can have a serious impact on mate selection and successful reproduction. Several species of Ornithoptera have been known to create hybrids if they have no access to their own species. Troides oblongamaculatus females have been known to choose to mate with other species such as Ornithoptera priamus poseidon, which will attempt mating if their own species is not to be found near by. The females will typically resist mating attempts by covering their abdomen with their forewings or dropping to the ground, making mating near impossible. Although the females usually resist these mating attempts, they have been noted to be more susceptible if they have not had previous encounters with males of their own species. Male courtship Some male Ornithoptera species demonstrate courtship behaviour. Ornithoptera priamus posedion males will approach a female carefully, and examine the female for several minutes. After consideration, the male may choose to hover twenty to thirty centimeters above the female, displaying the bright yellow marking on its hindwings. Meanwhile, the forewings will move forward, exposing the abdomen and androconial hair tufts. Mating is only attempted when the female has ceased to flap her wings. After about thirty seconds of the display, the male will attempt copulation. Cryptic choice: sperm competition and postcopulatory guarding In many animals, females often mate with more than one male. Males who are able will adapt strategies such as postcopulatory guarding to ensure the paternity of the offspring. Following insemination, it is common for the male Ornithoptera to produce a mating plug, which will seal the ostium bursae and prevent remating by the female, as new sperm is unable to enter the opening. The plug does not impede oviposition and may stay in place for the duration of the female's life. Sexual dimorphism Sexual dimorphism is very prominent in Ornithoptera species, the males being black with brightly colored markings of blue, green, orange or yellow and the females are overall black or dark brown. The sexual dichromatism functions in mate recognition by the use of photoreceptors. Due to the protected nature of Ornithoptera it has been difficult to study the spectral sensitivities of the sexes although this difference in coloration alludes to the idea of sensory exploitation of the female's photoreceptors. The sensory bias of females to select for males with brighter wings has yet to be studied in Ornithoptera. Gyanandromorphism is a very rare condition in which an organism simultaneously expresses both male and female phenotypes. It is only observed in species that express strong sexual dimorphism. Gynandromorphs are suspected to be due to genetic errors associated with cell division such as nondisjunction, as well as fertilization of binucleate ova and fertilisation of multiple sperm that may fuse and act as a second nucleus. Ornithoptera is known to commonly exhibit this phenomenon, but little to no research has been successful in determining why. Those who experience this phenomenon, usually females, show male-pigmented tissues on their wings.	Biology and health sciences	Lepidoptera	Animals
2136925	https://en.wikipedia.org/wiki/Lip	Lip	The lips are a horizontal pair of soft appendages attached to the jaws and are the most visible part of the mouth of many animals, including humans. Vertebrate lips are soft, movable and serve to facilitate the ingestion of food (e.g. suckling and gulping) and the articulation of sound and speech. Human lips are also a somatosensory organ, and can be an erogenous zone when used in kissing and other acts of intimacy. Structure The upper and lower lips are referred to as the labium superius oris and labium inferius oris, respectively. The juncture where the lips meet the surrounding skin of the mouth area is the vermilion border, and the typically reddish area within the borders is called the vermilion zone. The vermilion border of the upper lip is known as the Cupid's bow. The fleshy protuberance located in the center of the upper lip is a tubercle known by various terms including the procheilon (also spelled prochilon), the "tuberculum labii superioris", and the "labial tubercle". The vertical groove extending from the procheilon to the nasal septum is called the philtrum. The skin of the lip, with three to five cellular layers, is very thin compared to typical face skin, which has up to 16 layers. With light skin color, the lip skin contains fewer melanocytes (cells which produce melanin pigment, which give skin its color). Because of this, the blood vessels appear through the skin of the lips, which leads to their notable red coloring. With darker skin color this effect is less prominent, as in this case the skin of the lips contains more melanin and thus is visually darker. The skin of the lip forms the border between the exterior skin of the face, and the interior mucous membrane of the inside of the mouth. The lip skin is not hairy and does not have sweat glands. Therefore, it does not have the usual protection layer of sweat and body oils which keep the skin smooth, inhibit pathogens, and regulate warmth. For these reasons, the lips dry out faster and become chapped more easily. The lower lip is formed from the mandibular prominence, a branch of the first pharyngeal arch. The lower lip covers the anterior body of the mandible. It is lowered by the depressor labii inferioris muscle and the orbicularis oris borders it inferiorly. The upper lip covers the anterior surface of the body of the maxilla. Its upper half is of usual skin color and has a depression at its center, directly under the nasal septum, called the philtrum, which is Latin for "lower nose", while its lower half is a markedly different, red-colored skin tone more similar to the color of the inside of the mouth, and the term vermillion refers to the colored portion of either the upper or lower lip. It is raised by the levator labii superioris and is connected to the lower lip by the thin lining of the lip itself. Thinning of the vermilion of the upper lip and flattening of the philtrum are two of the facial characteristics of fetal alcohol syndrome, a lifelong disability caused by the mother's consumption of alcohol during pregnancy. Microanatomy The skin of the lips is stratified squamous epithelium. The mucous membrane is represented by a large area in the sensory cortex, and is therefore highly sensitive. The frenulum labii inferioris is the frenulum of the lower lip. The frenulum labii superioris is the frenulum of the upper lip. Nerve supply Trigeminal nerve The infraorbital nerve is a branch of the maxillary branch. It supplies not only the upper lip but also much of the skin of the face between the upper lip and the lower eyelid, except for the bridge of the nose. The mental nerve is a branch of the mandibular branch (via the inferior alveolar nerve). It supplies the skin and mucous membrane of the lower lip and labial gingiva (gum) anteriorly. Blood supply The facial artery is one of the six non-terminal branches of the external carotid artery. This artery supplies both lips by its superior and inferior labial branches. Each of the two branches bifurcate and anastomose with their companion branch from the other terminal. Muscles The muscles acting on the lips are considered part of the muscles of facial expression. All muscles of facial expression are derived from the mesoderm of the second pharyngeal arch and are therefore supplied (motor supply) by the nerve of the second pharyngeal arch, the facial nerve (7th cranial nerve). The muscles of facial expression are all specialized members of the panniculus carnosus, which attach to the dermis and so wrinkle or dimple the overlying skin. Functionally, the muscles of facial expression are arranged in groups around the orbits, nose, and mouth. The muscles acting on the lips: Buccinator Orbicularis oris (a complex of muscles, formerly thought to be a single sphincter or ring of muscle) Anchor point for several muscles Modiolus Lip elevation Levator labii superioris levator labii superioris alaeque nasi Levator anguli oris Zygomaticus minor Zygomaticus major Lip depression Risorius Depressor anguli oris Depressor labii inferioris Mentalis Functions Food intake Because they have their own muscles and bordering muscles, the lips are easily movable. Lips are used for eating functions, like holding food or to get it in the mouth. In addition, lips serve to close the mouth airtight shut, to hold food and drink inside, and to keep out unwanted objects. Through making a narrow funnel with the lips, the suction of the mouth is increased. This suction is essential for babies to breast feed. Lips can also be used to suck in other contexts, such as sucking on a straw to drink liquids. Articulation The lips serve for creating different sounds—mainly labial, bilabial, and labiodental consonant sounds as well as vowel rounding—and thus are an important part of the speech apparatus. The lips enable whistling and the performing of wind instruments such as the trumpet, clarinet, flute, and saxophone. People who have hearing loss may unconsciously or consciously lip read to understand speech without needing to perceive the actual sounds, and visual cues from the lips affect the perception of what sounds have been heard, for example the McGurk effect. Tactile organ The lip has many nerve endings and reacts as part of the tactile (touch) senses. Lips are very sensitive to touch, warmth, and cold. It is therefore an important aid for exploring unknown objects for babies and toddlers. Erogenous zone Because of their high number of nerve endings, the lips are an erogenous zone. The lips therefore play a crucial role in kissing and other acts of intimacy. A woman's lips are also a visible expression of her fertility. In studies performed on the science of human attraction, psychologists have concluded that a woman's facial and sexual attractiveness is closely linked to the makeup of her hormones during puberty and development. Contrary to the effects of testosterone on a man's facial structure, the effects of a woman's oestrogen levels serve to maintain a relatively "childlike" and youthful facial structure during puberty and during final maturation. It has been shown that the more oestrogen a woman has, the larger her eyes and the fuller her lips, characteristics which are perceived as more feminine. Surveys performed by sexual psychologists have also found that universally, men find a woman's full lips to be more sexually attractive than lips that are less so. A woman's lips are therefore sexually attractive to males because they serve as a biological indicator of a woman's health and fertility. A woman's lipstick (or collagen lip enhancement) attempts to take advantage of this fact by creating the illusion that a woman has more oestrogen than she actually has and thus that she is more fertile and attractive. Lip size is linked to sexual attraction in both men and women. Women are attracted to men with masculine lips that are more middle size and not too big or too small; they are to be rugged and sensual. In general, the researchers found that a small nose, big eyes and voluptuous lips are sexually attractive both in men and women. The lips may temporarily swell during sexual arousal due to engorgement with blood. Facial expression The lips contribute substantially to facial expressions. The lips visibly express emotions such as a smile or frown, iconically by the curve of the lips forming an up-open or down-open arc, respectively. Lips can also be made pouty when whining or perky to be provocative. Open questions The function of the abrupt change in skin structure between the lips and surrounding face (in particular, the function of the less keratinized vermillion and the white roll) is not completely understood. Possible reasons for the difference may include advantages to somatosensory function, better communication of facial expressions, and/or emphasis of the lips' slight sexual dimorphism as a secondary sex characteristic. Clinical significance As an organ of the body, the lip can be a focus of disease or show symptoms of a disease: One of the most frequent changes of the lips is a blue coloring due to cyanosis; the blood contains less oxygen and thus has a dark red to blue color, which shows through the thin skin. Cyanosis is the reason why corpses sometimes have blue lips. In cold weather cyanosis can appear, so especially in the winter, blue lips may not be an uncommon sight. Inflammation of the lips is termed cheilitis. This can be in several forms such as chapped lips (dry, peeling lips), angular cheilitis (inflammation of the corners of the mouth), herpes labialis (cold sore, a form of herpes simplex) and actinic cheilitis (chronically sun damaged lips). Cleft lip is a type of birth defect that can be successfully treated with surgery. Carcinoma (a malignant cancer that arises from epithelial cells) at the lips is caused predominantly by using tobacco and overexposure of sunlight. Alcohol appears to increase the carcinoma risk associated with tobacco use. It is most often a diffuse and often hyperkeratinised lesion, occasionally has the form of nodules and grows infiltratively, and can also be a combination of the two types. It more often occurs at the lower lip, where it is also much more malign. Lower lip carcinoma is exclusively planocellular carcinoma, whereas at the upper lip, it can also be basocellular carcinoma. Society and culture Lips are often viewed as a symbol of sensuality and sexuality. This has many origins; above all, the lips are a very sensitive erogenous and tactile organ. Furthermore, in many cultures of the world, a woman's mouth and lips are veiled because of their representative association with the vulva, and because of their role as a woman's secondary sexual organ. As part of the mouth, the lips are also associated with the symbolism associated with the mouth as orifice by which food is taken in. The lips are also linked symbolically to neonatal psychology (see for example oral stage of the psychology according to Sigmund Freud). Lip piercing or lip augmentation is sometimes carried out for cosmetic reasons. Products designed for use on the lips include lipstick, lip gloss and lip balm. Other animals In most vertebrates, the lips are relatively unimportant folds of tissue lying just outside the jaws. However, in mammals, they become much more prominent, being separated from the jaws by a deep cleft (a notable exception being the naked mole-rat, whose lips close behind the front teeth). They are also more mobile in mammals than in other groups since it is only in this group that they have any attached muscles. In some teleost fish, the lips may be modified to carry sensitive barbels. In birds and turtles, the lips are hard and keratinous, forming a solid beak. Clevosaurids like Clevosaurus are notable for the presence of bone "lips"; in these species the tooth-like jaw projections common to all sphenodontians form a beak-like edge around the jaws, protecting the teeth within.	Biology and health sciences	Gastrointestinal tract	Biology
2137292	https://en.wikipedia.org/wiki/Drag%20%28physics%29	Drag (physics)	In fluid dynamics, drag, sometimes referred to as fluid resistance, is a force acting opposite to the relative motion of any object moving with respect to a surrounding fluid. This can exist between two fluid layers, two solid surfaces, or between a fluid and a solid surface. Drag forces tend to decrease fluid velocity relative to the solid object in the fluid's path. Unlike other resistive forces, drag force depends on velocity. Drag force is proportional to the relative velocity for low-speed flow and is proportional to the velocity squared for high-speed flow. This distinction between low and high-speed flow is measured by the Reynolds number. Examples Examples of drag include: Net aerodynamic or hydrodynamic force: Drag acting opposite to the direction of movement of a solid object such as cars, aircraft, and boat hulls. Viscous drag of fluid in a pipe: Drag force on the immobile pipe decreases fluid velocity relative to the pipe. In the physics of sports, drag force is necessary to explain the motion of balls, javelins, arrows, and frisbees and the performance of runners and swimmers. For a top sprinter, overcoming drag can require 5% of their energy output. Types Types of drag are generally divided into the following categories: form drag due to the size and shape of a body skin friction drag or viscous drag due to the friction between the fluid and a surface which may be the outside of an object, or inside such as the bore of a pipe The effect of streamlining on the relative proportions of skin friction and form drag is shown for two different body sections: An airfoil, which is a streamlined body, and a cylinder, which is a bluff body. Also shown is a flat plate illustrating the effect that orientation has on the relative proportions of skin friction, and pressure difference between front and back. A body is known as bluff or blunt when the source of drag is dominated by pressure forces, and streamlined if the drag is dominated by viscous forces. For example, road vehicles are bluff bodies. For aircraft, pressure and friction drag are included in the definition of parasitic drag. Parasite drag is often expressed in terms of a hypothetical. Parasitic drag This is the area of a flat plate perpendicular to the flow. It is used when comparing the drag of different aircraft For example, the Douglas DC-3 has an equivalent parasite area of and the McDonnell Douglas DC-9, with 30 years of advancement in aircraft design, an area of although it carried five times as many passengers. lift-induced drag appears with wings or a lifting body in aviation and with semi-planing or planing hulls for watercraft wave drag (aerodynamics) is caused by the presence of shockwaves and first appears at subsonic aircraft speeds when local flow velocities become supersonic. The wave drag of the supersonic Concorde prototype aircraft was reduced at Mach 2 by 1.8% by applying the area rule which extended the rear fuselage on the production aircraft. wave resistance (ship hydrodynamics) or wave drag occurs when a solid object is moving along a fluid boundary and making surface waves boat-tail drag on an aircraft is caused by the angle with which the rear fuselage, or engine nacelle, narrows to the engine exhaust diameter. Lift-induced drag and parasitic drag Lift-induced drag Lift-induced drag (also called induced drag) is drag which occurs as the result of the creation of lift on a three-dimensional lifting body, such as the wing or propeller of an airplane. Induced drag consists primarily of two components: drag due to the creation of trailing vortices (vortex drag); and the presence of additional viscous drag (lift-induced viscous drag) that is not present when lift is zero. The trailing vortices in the flow-field, present in the wake of a lifting body, derive from the turbulent mixing of air from above and below the body which flows in slightly different directions as a consequence of creation of lift. With other parameters remaining the same, as the lift generated by a body increases, so does the lift-induced drag. This means that as the wing's angle of attack increases (up to a maximum called the stalling angle), the lift coefficient also increases, and so too does the lift-induced drag. At the onset of stall, lift is abruptly decreased, as is lift-induced drag, but viscous pressure drag, a component of parasite drag, increases due to the formation of turbulent unattached flow in the wake behind the body. Parasitic drag Parasitic drag, or profile drag, is drag caused by moving a solid object through a fluid. Parasitic drag is made up of multiple components including viscous pressure drag (form drag), and drag due to surface roughness (skin friction drag). Additionally, the presence of multiple bodies in relative proximity may incur so called interference drag, which is sometimes described as a component of parasitic drag. In aviation, induced drag tends to be greater at lower speeds because a high angle of attack is required to maintain lift, creating more drag. However, as speed increases the angle of attack can be reduced and the induced drag decreases. Parasitic drag, however, increases because the fluid is flowing more quickly around protruding objects increasing friction or drag. At even higher speeds (transonic), wave drag enters the picture. Each of these forms of drag changes in proportion to the others based on speed. The combined overall drag curve therefore shows a minimum at some airspeed - an aircraft flying at this speed will be at or close to its optimal efficiency. Pilots will use this speed to maximize endurance (minimum fuel consumption), or maximize gliding range in the event of an engine failure. The drag equation Drag depends on the properties of the fluid and on the size, shape, and speed of the object. One way to express this is by means of the drag equation: where is the drag force, is the density of the fluid, is the speed of the object relative to the fluid, is the cross sectional area, and is the drag coefficient – a dimensionless number. The drag coefficient depends on the shape of the object and on the Reynolds number where is some characteristic diameter or linear dimension. Actually, is the equivalent diameter of the object. For a sphere, is the D of the sphere itself. For a rectangular shape cross-section in the motion direction, , where a and b are the rectangle edges. is the kinematic viscosity of the fluid (equal to the dynamic viscosity divided by the density ). At low , is asymptotically proportional to , which means that the drag is linearly proportional to the speed, i.e. the drag force on a small sphere moving through a viscous fluid is given by the Stokes Law: At high , is more or less constant, but drag will vary as the square of the speed varies. The graph to the right shows how varies with for the case of a sphere. Since the power needed to overcome the drag force is the product of the force times speed, the power needed to overcome drag will vary as the square of the speed at low Reynolds numbers, and as the cube of the speed at high numbers. It can be demonstrated that drag force can be expressed as a function of a dimensionless number, which is dimensionally identical to the Bejan number. Consequently, drag force and drag coefficient can be a function of Bejan number. In fact, from the expression of drag force it has been obtained: and consequently allows expressing the drag coefficient as a function of Bejan number and the ratio between wet area and front area : where is the Reynolds number related to fluid path length L. At high velocity As mentioned, the drag equation with a constant drag coefficient gives the force moving through fluid a relatively large velocity, i.e. high Reynolds number, Re > ~1000. This is also called quadratic drag. The derivation of this equation is presented at . The reference area A is often the orthographic projection of the object, or the frontal area, on a plane perpendicular to the direction of motion. For objects with a simple shape, such as a sphere, this is the cross sectional area. Sometimes a body is a composite of different parts, each with a different reference area (drag coefficient corresponding to each of those different areas must be determined). In the case of a wing, the reference areas are the same, and the drag force is in the same ratio as the lift force. Therefore, the reference for a wing is often the lifting area, sometimes referred to as "wing area" rather than the frontal area. For an object with a smooth surface, and non-fixed separation points (like a sphere or circular cylinder), the drag coefficient may vary with Reynolds number Re, up to extremely high values (Re of the order 107). For an object with well-defined fixed separation points, like a circular disk with its plane normal to the flow direction, the drag coefficient is constant for Re > 3,500. The further the drag coefficient Cd is, in general, a function of the orientation of the flow with respect to the object (apart from symmetrical objects like a sphere). Power Under the assumption that the fluid is not moving relative to the currently used reference system, the power required to overcome the aerodynamic drag is given by: The power needed to push an object through a fluid increases as the cube of the velocity increases. For example, a car cruising on a highway at may require only to overcome aerodynamic drag, but that same car at requires . With a doubling of speeds, the drag/force quadruples per the formula. Exerting 4 times the force over a fixed distance produces 4 times as much work. At twice the speed, the work (resulting in displacement over a fixed distance) is done twice as fast. Since power is the rate of doing work, 4 times the work done in half the time requires 8 times the power. When the fluid is moving relative to the reference system, for example, a car driving into headwind, the power required to overcome the aerodynamic drag is given by the following formula: Where is the wind speed and is the object speed (both relative to ground). Velocity of a falling object Velocity as a function of time for an object falling through a non-dense medium, and released at zero relative-velocity v = 0 at time t = 0, is roughly given by a function involving a hyperbolic tangent (tanh): The hyperbolic tangent has a limit value of one, for large time t. In other words, velocity asymptotically approaches a maximum value called the terminal velocity vt: For an object falling and released at relative-velocity v = vi at time t = 0, with vi < vt, is also defined in terms of the hyperbolic tangent function: For vi > vt, the velocity function is defined in terms of the hyperbolic cotangent function: The hyperbolic cotangent also has a limit value of one, for large time t. Velocity asymptotically tends to the terminal velocity vt, strictly from above vt. For vi = vt, the velocity is constant: These functions are defined by the solution of the following differential equation: Or, more generically (where F(v) are the forces acting on the object beyond drag): For a potato-shaped object of average diameter d and of density ρobj, terminal velocity is about For objects of water-like density (raindrops, hail, live objects—mammals, birds, insects, etc.) falling in air near Earth's surface at sea level, the terminal velocity is roughly equal to with d in metre and vt in m/s. For example, for a human body ( ≈0.6 m) ≈70 m/s, for a small animal like a cat ( ≈0.2 m) ≈40 m/s, for a small bird ( ≈0.05 m) ≈20 m/s, for an insect ( ≈0.01 m) ≈9 m/s, and so on. Terminal velocity for very small objects (pollen, etc.) at low Reynolds numbers is determined by Stokes law. In short, terminal velocity is higher for larger creatures, and thus potentially more deadly. A creature such as a mouse falling at its terminal velocity is much more likely to survive impact with the ground than a human falling at its terminal velocity. Low Reynolds numbers: Stokes' drag The equation for viscous resistance or linear drag is appropriate for objects or particles moving through a fluid at relatively slow speeds (assuming there is no turbulence). Purely laminar flow only exists up to Re = 0.1 under this definition. In this case, the force of drag is approximately proportional to velocity. The equation for viscous resistance is: where: is a constant that depends on both the material properties of the object and fluid, as well as the geometry of the object; and is the velocity of the object. When an object falls from rest, its velocity will be where: is the density of the object, is density of the fluid, is the volume of the object, is the acceleration due to gravity (i.e., 9.8 m/s), and is mass of the object. The velocity asymptotically approaches the terminal velocity . For a given , denser objects fall more quickly. For the special case of small spherical objects moving slowly through a viscous fluid (and thus at small Reynolds number), George Gabriel Stokes derived an expression for the drag constant: where is the Stokes radius of the particle, and is the fluid viscosity. The resulting expression for the drag is known as Stokes' drag: For example, consider a small sphere with radius = 0.5 micrometre (diameter = 1.0 μm) moving through water at a velocity of 10 μm/s. Using 10−3 Pa·s as the dynamic viscosity of water in SI units, we find a drag force of 0.09 pN. This is about the drag force that a bacterium experiences as it swims through water. The drag coefficient of a sphere can be determined for the general case of a laminar flow with Reynolds numbers less than using the following formula: For Reynolds numbers less than 1, Stokes' law applies and the drag coefficient approaches ! Aerodynamics In aerodynamics, aerodynamic drag, also known as air resistance, is the fluid drag force that acts on any moving solid body in the direction of the air's freestream flow. From the body's perspective (near-field approach), the drag results from forces due to pressure distributions over the body surface, symbolized . Forces due to skin friction, which is a result of viscosity, denoted . Alternatively, calculated from the flow field perspective (far-field approach), the drag force results from three natural phenomena: shock waves, vortex sheet, and viscosity. Overview of aerodynamics When the airplane produces lift, another drag component results. Induced drag, symbolized , is due to a modification of the pressure distribution due to the trailing vortex system that accompanies the lift production. An alternative perspective on lift and drag is gained from considering the change of momentum of the airflow. The wing intercepts the airflow and forces the flow to move downward. This results in an equal and opposite force acting upward on the wing which is the lift force. The change of momentum of the airflow downward results in a reduction of the rearward momentum of the flow which is the result of a force acting forward on the airflow and applied by the wing to the air flow; an equal but opposite force acts on the wing rearward which is the induced drag. Another drag component, namely wave drag, , results from shock waves in transonic and supersonic flight speeds. The shock waves induce changes in the boundary layer and pressure distribution over the body surface. Therefore, there are three ways of categorizing drag. Pressure drag and friction drag Profile drag and induced drag Vortex drag, wave drag and wake drag The pressure distribution acting on a body's surface exerts normal forces on the body. Those forces can be added together and the component of that force that acts downstream represents the drag force, . The nature of these normal forces combines shock wave effects, vortex system generation effects, and wake viscous mechanisms. Viscosity of the fluid has a major effect on drag. In the absence of viscosity, the pressure forces acting to hinder the vehicle are canceled by a pressure force further aft that acts to push the vehicle forward; this is called pressure recovery and the result is that the drag is zero. That is to say, the work the body does on the airflow is reversible and is recovered as there are no frictional effects to convert the flow energy into heat. Pressure recovery acts even in the case of viscous flow. Viscosity, however results in pressure drag and it is the dominant component of drag in the case of vehicles with regions of separated flow, in which the pressure recovery is infective. The friction drag force, which is a tangential force on the aircraft surface, depends substantially on boundary layer configuration and viscosity. The net friction drag, , is calculated as the downstream projection of the viscous forces evaluated over the body's surface. The sum of friction drag and pressure (form) drag is called viscous drag. This drag component is due to viscosity. History The idea that a moving body passing through air or another fluid encounters resistance had been known since the time of Aristotle. According to Mervyn O'Gorman, this was named "drag" by Archibald Reith Low. Louis Charles Breguet's paper of 1922 began efforts to reduce drag by streamlining. Breguet went on to put his ideas into practice by designing several record-breaking aircraft in the 1920s and 1930s. Ludwig Prandtl's boundary layer theory in the 1920s provided the impetus to minimise skin friction. A further major call for streamlining was made by Sir Melvill Jones who provided the theoretical concepts to demonstrate emphatically the importance of streamlining in aircraft design. In 1929 his paper 'The Streamline Airplane' presented to the Royal Aeronautical Society was seminal. He proposed an ideal aircraft that would have minimal drag which led to the concepts of a 'clean' monoplane and retractable undercarriage. The aspect of Jones's paper that most shocked the designers of the time was his plot of the horse power required versus velocity, for an actual and an ideal plane. By looking at a data point for a given aircraft and extrapolating it horizontally to the ideal curve, the velocity gain for the same power can be seen. When Jones finished his presentation, a member of the audience described the results as being of the same level of importance as the Carnot cycle in thermodynamics. Power curve in aviation The interaction of parasitic and induced drag vs. airspeed can be plotted as a characteristic curve, illustrated here. In aviation, this is often referred to as the power curve, and is important to pilots because it shows that, below a certain airspeed, maintaining airspeed counterintuitively requires more thrust as speed decreases, rather than less. The consequences of being "behind the curve" in flight are important and are taught as part of pilot training. At the subsonic airspeeds where the "U" shape of this curve is significant, wave drag has not yet become a factor, and so it is not shown in the curve. Wave drag in transonic and supersonic flow Wave drag, sometimes referred to as compressibility drag, is drag that is created when a body moves in a compressible fluid and at the speed that is close to the speed of sound in that fluid. In aerodynamics, wave drag consists of multiple components depending on the speed regime of the flight. In transonic flight, wave drag is the result of the formation of shockwaves in the fluid, formed when local areas of supersonic (Mach number greater than 1.0) flow are created. In practice, supersonic flow occurs on bodies traveling well below the speed of sound, as the local speed of air increases as it accelerates over the body to speeds above Mach 1.0. However, full supersonic flow over the vehicle will not develop until well past Mach 1.0. Aircraft flying at transonic speed often incur wave drag through the normal course of operation. In transonic flight, wave drag is commonly referred to as transonic compressibility drag. Transonic compressibility drag increases significantly as the speed of flight increases towards Mach 1.0, dominating other forms of drag at those speeds. In supersonic flight (Mach numbers greater than 1.0), wave drag is the result of shockwaves present in the fluid and attached to the body, typically oblique shockwaves formed at the leading and trailing edges of the body. In highly supersonic flows, or in bodies with turning angles sufficiently large, unattached shockwaves, or bow waves will instead form. Additionally, local areas of transonic flow behind the initial shockwave may occur at lower supersonic speeds, and can lead to the development of additional, smaller shockwaves present on the surfaces of other lifting bodies, similar to those found in transonic flows. In supersonic flow regimes, wave drag is commonly separated into two components, supersonic lift-dependent wave drag and supersonic volume-dependent wave drag. The closed form solution for the minimum wave drag of a body of revolution with a fixed length was found by Sears and Haack, and is known as the Sears-Haack Distribution. Similarly, for a fixed volume, the shape for minimum wave drag is the Von Karman Ogive. The Busemann biplane theoretical concept is not subject to wave drag when operated at its design speed, but is incapable of generating lift in this condition. d'Alembert's paradox In 1752 d'Alembert proved that potential flow, the 18th century state-of-the-art inviscid flow theory amenable to mathematical solutions, resulted in the prediction of zero drag. This was in contradiction with experimental evidence, and became known as d'Alembert's paradox. In the 19th century the Navier–Stokes equations for the description of viscous flow were developed by Saint-Venant, Navier and Stokes. Stokes derived the drag around a sphere at very low Reynolds numbers, the result of which is called Stokes' law. In the limit of high Reynolds numbers, the Navier–Stokes equations approach the inviscid Euler equations, of which the potential-flow solutions considered by d'Alembert are solutions. However, all experiments at high Reynolds numbers showed there is drag. Attempts to construct inviscid steady flow solutions to the Euler equations, other than the potential flow solutions, did not result in realistic results. The notion of boundary layers—introduced by Prandtl in 1904, founded on both theory and experiments—explained the causes of drag at high Reynolds numbers. The boundary layer is the thin layer of fluid close to the object's boundary, where viscous effects remain important even when the viscosity is very small (or equivalently the Reynolds number is very large).	Physical sciences	Fluid mechanics	null
15845253	https://en.wikipedia.org/wiki/Viral%20disease	Viral disease	A viral disease (or viral infection) occurs when an organism's body is invaded by pathogenic viruses, and infectious virus particles (virions) attach to and enter susceptible cells. Examples are the common cold, gastroenteritis,corona,flu,pneumonia. Structural characteristics Basic structural characteristics, such as genome type, virion shape and replication site, generally share the same features among virus species within the same family. Double-stranded DNA families: three are non-enveloped (Adenoviridae, Papillomaviridae and Polyomaviridae) and two are enveloped (Herpesviridae and Poxviridae). All of the non-enveloped families have icosahedral capsids. Partly double-stranded DNA viruses: Hepadnaviridae. These viruses are enveloped. One family of single-stranded DNA viruses infects humans: Parvoviridae. These viruses are non-enveloped. Positive single-stranded RNA families: three non-enveloped (Astroviridae, Caliciviridae and Picornaviridae) and four enveloped (Coronaviridae, Flaviviridae, Retroviridae and Togaviridae). All the non-enveloped families have icosahedral nucleocapsids. Negative single-stranded RNA families: Arenaviridae, Bunyaviridae, Filoviridae, Orthomyxoviridae, Paramyxoviridae and Rhabdoviridae. All are enveloped with helical nucleocapsids. Double-stranded RNA genome: Reoviridae. The Hepatitis D virus has not yet been assigned to a family, but is clearly distinct from the other families infecting humans. Viruses known to infect humans that have not been associated with disease: the family Anelloviridae and the genus Dependovirus. Both of these taxa are non-enveloped single-stranded DNA viruses. Pragmatic rules Human-infecting virus families offer rules that may assist physicians and medical microbiologists/virologists. As a general rule, DNA viruses replicate within the cell nucleus while RNA viruses replicate within the cytoplasm. Exceptions are known to this rule: poxviruses replicate within the cytoplasm and orthomyxoviruses and hepatitis D virus (RNA viruses) replicate within the nucleus. Segmented genomes: Bunyaviridae, Orthomyxoviridae, Arenaviridae, and Reoviridae (acronym BOAR). All are RNA viruses. Viruses transmitted almost exclusively by arthropods: Bunyavirus, Flavivirus, and Togavirus. Some Reoviruses are transmitted from arthropod vectors. All are RNA viruses. One family of enveloped viruses causes gastroenteritis (Coronaviridae). All other viruses associated with gastroenteritis are non-enveloped. Baltimore group This group of analysts defined multiple categories of virus. Groups: I - dsDNA II - ssDNA III - dsRNA IV - positive-sense ssRNA V - negative-sense ssRNA VI - ssRNA-RT VII - dsDNA-RT Clinical characteristics The clinical characteristics of viruses may differ substantially among species within the same family:	Biology and health sciences	Concepts	Health
958996	https://en.wikipedia.org/wiki/Messier%2015	Messier 15	Messier 15 or M15 (also designated NGC 7078 and sometimes known as the Great Pegasus Cluster) is a globular cluster in the constellation Pegasus. It was discovered by Jean-Dominique Maraldi in 1746 and included in Charles Messier's catalogue of comet-like objects in 1764. At an estimated billion years old, it is one of the oldest known globular clusters. Characteristics M 15 is about 35,700 light-years from Earth, and 175 light-years in diameter. It has an absolute magnitude of −9.2, which translates to a total luminosity of 360,000 times that of the Sun. Messier 15 is one of the most densely packed globulars known in the Milky Way galaxy. Its core has undergone a contraction known as "core collapse" and it has a central density cusp with an enormous number of stars surrounding what may be a central black hole. Home to over 100,000 stars, the cluster is notable for containing a large number of variable stars (112) and pulsars (8), including one double neutron star system, M15-C. It also contains Pease 1, the first planetary nebula discovered within a globular cluster in 1928. Just three others have been found in globular clusters since then. Amateur astronomy At magnitude 6.2, M15 approaches naked eye visibility under good conditions and can be observed with binoculars or a small telescope, appearing as a fuzzy star. Telescopes with a larger aperture (at least 6 in. (150 mm)) will start to reveal individual stars, the brightest of which are of magnitude +12.6. The cluster appears 18 arc minutes in size (three tenths of a degree across). M15 is around 4° WNW of the brightest star of Pegasus, Epsilon Pegasi. X-ray sources Earth-orbiting satellites Uhuru and Chandra X-ray Observatory have detected two bright X-ray sources in this cluster: Messier 15 X-1 (4U 2129+12) and Messier 15 X-2. The former appears to be the first astronomical X-ray source detected in Pegasus. Gallery	Physical sciences	Notable star clusters	Astronomy
959018	https://en.wikipedia.org/wiki/Omega%20Nebula	Omega Nebula	The Omega Nebula is an H II region in the constellation Sagittarius. It was discovered by Philippe Loys de Chéseaux in 1745. Charles Messier catalogued it in 1764. It is by some of the richest starfields of the Milky Way, figuring in the northern two-thirds of Sagittarius. This feature is also known as the Swan Nebula, Checkmark Nebula, Lobster Nebula, and the Horseshoe Nebula, and catalogued as Messier 17 or M17 or NGC 6618. Characteristics The Omega Nebula is between 5,000 and 6,000 light-years from Earth and it spans some 15 light-years in diameter. The cloud of interstellar matter of which this nebula is a part is roughly 40 light-years in diameter and has a mass of 30,000 solar masses. The total mass of the Omega Nebula is an estimated 800 solar masses. It is considered one of the brightest and most massive star-forming regions of our galaxy. Its local geometry is similar to the Orion Nebula except that it is viewed edge-on rather than face-on. The open cluster NGC 6618 lies embedded in the nebulosity and causes the gases of the nebula to shine due to radiation from these hot, young stars; however, the actual number of stars in the nebula is much higher – up to 800, 100 of spectral type earlier than B9, and 9 of spectral type O, plus over a thousand stars in formation on its outer regions. It is also one of the youngest clusters known, with an age of just 1 million years. The luminous blue variable HD 168607, in the south-east part of the nebula, is generally assumed to be associated with it; its close neighbor, the blue hypergiant HD 168625, may be too. The Swan portion of M17, the Omega Nebula in the Sagittarius nebulosity is said to resemble a barber's pole. Early research The first attempt to accurately draw the nebula (as part of a series of sketches of nebulae) was made by John Herschel in 1833, and published in 1836. He described the nebula as such: A second, more detailed sketch was made during his visit to South Africa in 1837. The nebula was also studied by Johann von Lamont and separately by an undergraduate at Yale College, Mr Mason, starting from around 1836. When Herschel published his 1837 sketch in 1847, he wrote: Sketches were also made by William Lassell in 1862 using his four-foot telescope at Malta, and by M. Trouvelot from Cambridge, Massachusetts, and Edward Singleton Holden in 1875 using the twenty-six inch Clark refractor at the United States Naval Observatory. Observations by SOFIA In January 2020, the Stratospheric Observatory for Infrared Astronomy (SOFIA) provided new insights into the Omega Nebula. SOFIA's composite image revealed that blue areas (20 microns) near the center indicate gas heated by massive stars, while green areas (37 microns) trace dust warmed by massive stars and newborn stars. Nine previously unseen protostars were discovered primarily in the southern regions. Red areas near the edges represent cold dust detected by the Herschel Space Telescope (70 microns), and the white star field was observed by the Spitzer Space Telescope (3.6 microns). These observations suggest that parts of the nebula formed separately, contributing to its distinctive swan-like shape. Gallery	Physical sciences	Notable nebulae	Astronomy
959491	https://en.wikipedia.org/wiki/Hypovolemic%20shock	Hypovolemic shock	Hypovolemic shock is a form of shock caused by severe hypovolemia (insufficient blood volume or extracellular fluid in the body). It can be caused by severe dehydration or blood loss. Hypovolemic shock is a medical emergency; if left untreated, the insufficient blood flow can cause damage to organs, leading to multiple organ failure. In treating hypovolemic shock, it is important to determine the cause of the underlying hypovolemia, which may be the result of bleeding or other fluid losses. To minimize ischemic damage to tissues, treatment involves quickly replacing lost blood or fluids, with consideration of both rate and the type of fluids used. Tachycardia, a fast heart rate, is typically the first abnormal vital sign. When resulting from blood loss, trauma is the most common root cause, but severe blood loss can also happen in various body systems without clear traumatic injury. The body in hypovolemic shock prioritizes getting oxygen to the brain and heart, which reduces blood flow to nonvital organs and extremities, causing them to grow cold, look mottled, and exhibit delayed capillary refill. The lack of adequate oxygen delivery ultimately leads to a worsening increase in the acidity of the blood (acidosis). The "lethal triad" of ways trauma can lead to death is acidosis, hypothermia, and coagulopathy. It is possible for trauma to cause clotting problems even without resuscitation efforts. Damage control resuscitation is based on three principles: permissive hypotension: tries to balance temporary suboptimal perfusion to organs with conditions for halting blood loss by setting a goal of 90 mmHg systolic blood pressure hemostatic resuscitation: restoring blood volume in ways (with whole blood or equivalent) that interfere minimally with the natural process of stopping bleeding. damage control surgery. Signs and symptoms Symptoms of hypovolemic shock can be related to volume depletion, electrolyte imbalances, or acid–base disorders that accompany hypovolemic shock. Patients with volume depletion may complain of thirst, muscle cramps, and/or orthostatic hypotension. Severe hypovolemic shock can result in mesenteric and coronary ischemia that can cause abdominal or chest pain. Agitation, lethargy, or confusion may characterize brain mal-perfusion. Dry mucous membranes, decreased skin turgor, low jugular venous distention, tachycardia, and hypotension can be seen along with decreased urinary output. Patients in shock can appear cold, clammy, and cyanotic. Early signs and symptoms include tachycardia given rise to by catecholamine release; skin pallor due to vasoconstriction triggered by catecholamine release; hypotension followed by hypovolaemia and perhaps arising after myocardial insufficiency; and confusion, aggression, drowsiness and coma caused by cerebral hypoxia or acidosis. Tachypnoea owing to hypoxia and acidosis, general weakness caused by hypoxia and acidosis, thirst induced by hypovolaemia, and oliguria caused by reduced perfusion may also arise. Abnormal growing central venous pressure indicates either hypotension or hypovolemia. Tachycardia accompanied by declined urine outflow implies either tension pneumothorax, cardiac tamponade or cardiac failure which is thought secondary to cardiac contusion or ischaemic heart disease. Echocardiography in such case may be helpful to distinguish cardiac failure from other diseases. Cardiac failure manifests a weak contractibility myocardium; treatment with an inotropic drug such as dobutamine may be appropriate. Cause The annual incidence of shock of any etiology is 0.3 to 0.7 per 1000, with hemorrhagic shock being most common in the intensive care unit. Hypovolemic shock is the most common type of shock in children, most commonly due to diarrheal illness in the developing world. Hypovolemic shock occurs as a result of either blood loss or extracellular fluid loss. Blood loss Hemorrhagic shock is hypovolemic shock from blood loss. Traumatic injury is by far the most common cause of hemorrhagic shock, particularly blunt and penetrating trauma, followed by upper and lower gastrointestinal sources, such as gastrointestinal (GI) bleed. Other causes of hemorrhagic shock include bleed from an ectopic pregnancy, bleeding from surgical intervention, vaginal bleeding, and splenic rupture. Obstetrical, vascular, iatrogenic, and even urological sources have all been described. Bleeding may be either external or internal. A substantial amount of blood loss to the point of hemodynamic compromise may occur in the chest, abdomen, or the retroperitoneum. The thigh itself can hold up to 1 L to 2 L of blood. Localizing and controlling the source of bleeding is of utmost importance to the treatment of hemorrhagic shock. The sequence of the most-commonly-seen causes that lead to hemorrhagic type of hypovolemic shock is given in order of frequencies: blunt or penetrating trauma including multiple fractures absent from vessel impairment, upper gastrointestinal bleeding e.g., variceal hemorrhage, peptic ulcer., or lower GI bleeding e.g., diverticular, and arteriovenous malformation. Except for the two most common causes, the less common causes are intra-operative and post-operative bleeding, abdominal aortic rupture or left ventricle aneurysm rupture, aortic–enteric fistula, hemorrhagic pancreatitis, iatrogenic e.g., inadvertent biopsy of arteriovenous malformation, severed artery., tumors or abscess erosion into major vessels, post-partum hemorrhage, uterine or vaginal hemorrhage owing to infection, tumors, lacerations, spontaneous peritoneal hemorrhage caused by bleeding diathesis, and ruptured hematoma. Fluid loss In spite of hemorrhage, the amount of circulating blood in the body may drop as well when one loses excessive body fluid owing to non-hemorrhagic reasons. Hypovolemic shock as a result of extracellular fluid loss can be of the 4 etiologies. Gastrointestinal Gastrointestinal (GI) losses can occur via many different etiologies. The gastrointestinal tract usually secretes between 3 and 6 liters of fluid per day. However, most of this fluid is reabsorbed as only 100 to 200 mL are lost in the stool. Volume depletion occurs when the fluid ordinarily secreted by the GI tract cannot be reabsorbed. This occurs when there is retractable vomiting, diarrhea, or external drainage via stoma or fistulas. Kidneys Renal losses of salt and fluid can lead to hypovolemic shock. The kidneys usually excrete sodium and water in a manner that matches sodium intake and water intake. Diuretic therapy and osmotic diuresis from hyperglycemia can lead to excessive renal sodium and volume loss. In addition, there are several tubular and interstitial diseases beyond the scope of this article that cause severe salt-wasting nephropathy. Skin Fluid loss also can occur from the skin. In a hot and dry climate, skin fluid losses can be as high as 1 to 2 liters/hour. Patients with a skin barrier interrupted by burns or other skin lesions also can experience large fluid losses that lead to hypovolemic shock. Third-spacing Sequestration of fluid into a third space also can lead to volume loss and hypovolemic shock. Third-spacing of fluid can occur in intestinal obstruction, pancreatitis, obstruction of a major venous system, vascular endothelium or any other pathological condition that results in a massive inflammatory response. Pathophysiology Blood loss Hemorrhagic shock is due to the depletion of intravascular volume through blood loss to the point of being unable to match the tissues' demand for oxygen. As a result, mitochondria are no longer able to sustain aerobic metabolism for the production of oxygen and switch to the less efficient anaerobic metabolism to meet the cellular demand for adenosine triphosphate. In the latter process, pyruvate is produced and converted to lactic acid to regenerate nicotinamide adenine dinucleotide (NAD+) to maintain some degree of cellular respiration in the absence of oxygen. The body compensates for volume loss by increasing heart rate and contractility, followed by baroreceptor activation resulting in sympathetic nervous system activation and peripheral vasoconstriction. Typically, there is a slight increase in the diastolic blood pressure with narrowing of the pulse pressure. As diastolic ventricular filling continues to decline and cardiac output decreases, systolic blood pressure drops. Due to sympathetic nervous system activation, blood is diverted away from noncritical organs and tissues to preserve blood supply to vital organs such as the heart and brain. While prolonging heart and brain function, this also leads to other tissues being further deprived of oxygen causing more lactic acid production and worsening acidosis. This worsening acidosis along with hypoxemia, if left uncorrected, eventually causes the loss of peripheral vasoconstriction, worsening hemodynamic compromise, and death. The body's compensation varies by cardiopulmonary comorbidities, age, and vasoactive medications. Due to these factors, heart rate and blood pressure responses are extremely variable and, therefore, cannot be relied upon as the sole means of diagnosis. A key factor in the pathophysiology of hemorrhagic shock is the development of trauma-induced coagulopathy. Coagulopathy develops as a combination of several processes. The simultaneous loss of coagulation factors via hemorrhage, hemodilution with resuscitation fluids, and coagulation cascade dysfunction secondary to acidosis and hypothermia have been traditionally thought to be the cause of coagulopathy in trauma. However, this traditional model of trauma-induced coagulopathy may be too limited. Further studies have shown that a degree of coagulopathy begins in 25% to 56% of patients before initiation of the resuscitation. This has led to the recognition of trauma-induced coagulopathy as the sum of two distinct processes: acute coagulopathy of trauma and resuscitation-induced coagulopathy. Trauma-induced coagulopathy is acutely worsened by the presence of acidosis and hypothermia. The activity of coagulation factors, fibrinogen depletion, and platelet quantity are all adversely affected by acidosis. Hypothermia (less than 34 C) compounds coagulopathy by impairing coagulation and is an independent risk factor for death in hemorrhagic shock. Fluid loss Hypovolemic shock results from depletion of intravascular volume, whether by extracellular fluid loss or blood loss. The body compensates with increased sympathetic tone resulting in increased heart rate, increased cardiac contractility, and peripheral vasoconstriction. The first changes in vital signs seen in hypovolemic shock include an increase in diastolic blood pressure with narrowed pulse pressure. As volume status continues to decrease, systolic blood pressure drops. As a result, oxygen delivery to vital organs is unable to meet the oxygen needs of the cells. Cells switch from aerobic metabolism to anaerobic metabolism, resulting in lactic acidosis. As sympathetic drive increases, blood flow is diverted from other organs to preserve blood flow to the heart and brain. This propagates tissue ischemia and worsens lactic acidosis. If not corrected, there will be worsening hemodynamic compromise and, eventually, death. Diagnosis Shock index (SI) has been defined as ; SI≥0.6 is a clinical shock. Such ratio value is clinically employed to determine the scope or emergence of shock. The SI correlates with the extent of hypovolemia and thus may facilitate the early identification of severely injured patients threatened by complications due to blood loss and therefore need urgent treatment, i.e. blood transfusion. Data presented as n (%), mean ± standard deviation or median (interquartile range (IQR)). n = 21,853; P <0.001 for all parameters. ED Emergency department, GCS Glasgow coma scale, HR Heart rate, SBP Systolic blood pressure, SI = Shock index. Bleeding Recognizing the degree of blood loss via vital sign and mental status abnormalities is important. The American College of Surgeons Advanced Trauma Life Support (ATLS) hemorrhagic shock classification links the amount of blood loss to expected physiologic responses in a healthy 70 kg patient. As total circulating blood volume accounts for approximately 7% of total body weight, this equals approximately five liters in the average 70 kg male patient. Class 1: Volume loss up to 15% of total blood volume, approximately 750 mL. Heart rate is minimally elevated or normal. Typically, there is no change in blood pressure, pulse pressure, or respiratory rate. Class 2: Volume loss from 15% to 30% of total blood volume, from 750 mL to 1500 mL. Heart rate and respiratory rate become elevated (100 BPM to 120 BPM, 20 RR to 24 RR). Pulse pressure begins to narrow, but systolic blood pressure may be unchanged to slightly decreased. Class 3: Volume loss from 30% to 40% of total blood volume, from 1500 mL to 2000 mL. A significant drop in blood pressure and changes in mental status occur. Heart rate and respiratory rate are significantly elevated (more than 120 BPM). Urine output declines. Capillary refill is delayed. Class 4: Volume loss over 40% of total blood volume. Hypotension with narrow pulse pressure (less than 25 mmHg). Tachycardia becomes more pronounced (more than 120 BPM), and mental status becomes increasingly altered. Urine output is minimal or absent. Capillary refill is delayed. Again, the above is outlined for a healthy 70 kg individual. Clinical factors must be taken into account when assessing patients. For example, elderly patients taking beta blockers can alter the patient's physiologic response to decreased blood volume by inhibiting mechanism to increase heart rate. As another, patients with baseline hypertension may be functionally hypotensive with a systolic blood pressure of 110 mmHg. Non-bleeding Various laboratory values can be abnormal in hypovolemic shock. Patients can have increased BUN and serum creatinine as a result of pre-renal kidney failure. Hypernatremia or hyponatremia can result, as can hyperkalemia or hypokalemia. Lactic acidosis can result from increased anaerobic metabolism. However, the effect of acid–base balance can be variable as patients with large GI losses can become alkalotic. In cases of hemorrhagic shock, hematocrit and hemoglobin can be severely decreased. However, with a reduction in plasma volume, hematocrit and hemoglobin can be increased due to hemoconcentration. Low urinary sodium is commonly found in hypovolemic patients as the kidneys attempt to conserve sodium and water to expand the extracellular volume. However, sodium urine can be low in a euvolemic patient with heart failure, cirrhosis, or nephrotic syndrome. Fractional excretion of sodium under 1% is also suggestive of volume depletion. Elevated urine osmolality can also suggest hypovolemia. However, this number also can be elevated in the setting of impaired concentrating ability by the kidneys. Central venous pressure (CVP) is often used to assess volume status. However, its usefulness in determining volume responsiveness has recently come into question. Ventilator settings, chest wall compliance, and right-sided heart failure can compromise CVPs accuracy as a measure of volume status. Measurements of pulse pressure variation via various commercial devices has also been postulated as a measure of volume responsiveness. However, pulse pressure variation as a measure of fluid responsiveness is only valid in patients without spontaneous breaths or arrhythmias. The accuracy of pulse pressure variation also can be compromised in right heart failure, decreased lung or chest wall compliance, and high respiratory rates. Similar to examining pulse pressure variation, measuring respiratory variation in inferior vena cava diameter as a measure of volume responsiveness has only been validated in patients without spontaneous breaths or arrhythmias. Measuring the effect of passive leg raises on cardiac contractility by echo appears to be the most accurate measurement of volume responsiveness, although it is also subject to limitations. History and physical can often make the diagnosis of hypovolemic shock. For patients with hemorrhagic shock, a history of trauma or recent surgery is present. For hypovolemic shock due to fluid losses, history and physical should attempt to identify possible GI, renal, skin, or third-spacing as a cause of extracellular fluid loss. Although relatively nonsensitive and nonspecific, physical exam can be helpful in determining the presence of hypovolemic shock. Physical findings suggestive of volume depletion include dry mucous membranes, decreased skin turgor, and low jugular venous distention. Tachycardia and hypotension can be seen along with decreased urinary output. Differential diagnosis While hemorrhage is the most common cause of shock in the trauma patient, other causes of shock are to remain on the differential. Obstructive shock can occur in the setting of tension pneumothorax and cardiac tamponade. These etiologies should be uncovered in the primary survey. In the setting of head or neck trauma, an inadequate sympathetic response, or neurogenic shock, is a type of distributive shock that is caused by a decrease in peripheral vascular resistance. This is suggested by an inappropriately low heart rate in the setting of hypotension. Cardiac contusion and infarctions can result in cardiogenic shock. Finally, other causes should be considered that are not related to trauma or blood loss. In the undifferentiated patient with shock, septic shock and toxic causes are also on the differential. Management The first step in managing hemorrhagic shock is recognition. Ideally, This should occur before the development of hypotension. Close attention should be paid to physiological responses to low-blood volume. Tachycardia, tachypnea, and narrowing pulse pressure may be the initial signs. Cool extremities and delayed capillary refill are signs of peripheral vasoconstriction. Bleeding In the setting of trauma, an algorithmic approach via the primary and secondary surveys is suggested by ATLS. Physical exam and radiological evaluations can help localize sources of bleeding. A trauma ultrasound, or Focused Assessment with Sonography for Trauma (FAST), has been incorporated in many circumstances into the initial surveys. The specificity of a FAST scan has been reported above 99%, but a negative ultrasound does not rule out intra-abdominal pathology. With a broader understanding of the pathophysiology of hemorrhagic shock, treatment in trauma has expanded from a simple massive transfusion method to a more comprehensive management strategy of "damage control resuscitation". The concept of damage control resuscitation focuses on permissive hypotension, hemostatic resuscitation, and hemorrhage control to adequately treat the "lethal triad" of coagulopathy, acidosis, and hypothermia that occurs in trauma. Hypotensive resuscitation has been suggested for the hemorrhagic shock patient without head trauma. The aim is to achieve a systolic blood pressure of 90 mmHg in order to maintain tissue perfusion without inducing re-bleeding from recently clotted vessels. Permissive hypotension is a means of restricting fluid administration until hemorrhage is controlled while accepting a short period of suboptimal end-organ perfusion. Studies regarding permissive hypotension have yielded conflicting results and must take into account type of injury (penetrating versus blunt), the likelihood of intracranial injury, the severity of the injury, as well as proximity to a trauma center and definitive hemorrhage control. The quantity, type of fluids to be used, and endpoints of resuscitation remain topics of much study and debate. For crystalloid resuscitation, normal saline and lactated ringers are the most commonly used fluids. Normal saline has the drawback of causing a non-anion gap hyperchloremic metabolic acidosis due to the high chloride content, while lactated ringers can cause a metabolic alkalosis as lactate metabolism regenerates into bicarbonate. Recent trends in damage control resuscitation focus on "hemostatic resuscitation" which pushes for early use of blood products rather than an abundance of crystalloids in order to minimize the metabolic derangement, resuscitation-induced coagulopathy, and the hemodilution that occurs with crystalloid resuscitation. The end goal of resuscitation and the ratios of blood products remain at the center of much study and debate. A recent study has shown no significant difference in mortality at 24 hours or 30 days between ratios of 1:1:1 and 1:1:2 of plasma to platelets to packed RBCs. However, patients that received the more balanced ratio of 1:1:1 were less likely to die as a result of exsanguination in 24 hours and were more likely to achieve hemostasis. Additionally, reduction in time to first plasma transfusion has shown a significant reduction in mortality in damage control resuscitation. In addition to blood products, products that prevent the breakdown of fibrin in clots, or antifibrinolytics, have been studied for their utility in the treatment of hemorrhagic shock in the trauma patient. Several antifibrinolytics have been shown to be safe and effective in elective surgery. The CRASH-2 study was a randomized control trial of tranexamic acid versus placebo in trauma has been shown to decrease overall mortality when given in the first three hours of injury. Follow-up analysis shows additional benefit to tranexamic acid when given in the first three hours after surgery. Damage control resuscitation is to occur in conjunction with prompt intervention to control the source of bleeding. Strategies may differ depending on proximity to definitive treatment. For patients in hemorrhagic shock, early use of blood products over crystalloid resuscitation results in better outcomes. Balanced transfusion using 1:1:1 or 1:1:2 of plasma to platelets to packed red blood cells results in better hemostasis. Anti-fibrinolytic administration to patients with severe bleed within 3 hours of traumatic injury appears to decrease death from major bleed as shown in the CRASH-2 trial. Research on oxygen-carrying substitutes as an alternative to packed red blood cells is ongoing, although no blood substitutes have been approved for use in the United States. Fluid loss For patients in hypovolemic shock due to fluid losses, the exact fluid deficit cannot be determined. Therefore, it is prudent to start with 2 liters of isotonic crystalloid solution infused rapidly as an attempt to quickly restore tissue perfusion. Fluid repletion can be monitored by measuring blood pressure, urine output, mental status, and peripheral edema. Multiple modalities exist for measuring fluid responsiveness such as ultrasound, central venous pressure monitoring, and pulse pressure fluctuation as described above. Vasopressors may be used if blood pressure does not improve with fluids. Crystalloid fluid resuscitation is preferred over colloid solutions for severe volume depletion not due to bleeding. The type of crystalloid used to resuscitate the patient can be individualized based on the patients' chemistries, estimated volume of resuscitation, acid/base status, and physician or institutional preferences. Isotonic saline is hyperchloremic relative to blood plasma, and resuscitation with large amounts can lead to hyperchloremic metabolic acidosis. Several other isotonic fluids with lower chloride concentrations exist, such as lactated Ringer's solution or PlasmaLyte. These solutions are often referred to as buffered or balanced crystalloids. Some evidence suggests that patients who need large volume resuscitation may have a less renal injury with restrictive chloride strategies and use of balanced crystalloids. Crystalloid solutions are equally as effective and much less expensive than colloid. Commonly used colloid solutions include those containing albumin or hyperoncotic starch. Studies examining albumin solutions for resuscitation have not shown improved outcomes, while other studies have shown resuscitation with hyper-oncotic starch leads to increased mortality rate and renal failure. Patients in shock can appear cold, clammy, and cyanotic. Hypothermia increases the mortality rate of patients with hypovolemic shock. It is advised to keep the patient warm for the sake of maintaining the temperatures of all kinds of fluids inside the patient. Monitoring parameters Oxygen saturation by pulse oximetry (SpO2). Respiratory rate. Pulse rate. Arterial blood pressure. Pulse pressure. Central venous pressure. Urine output. Base deficit and/or lactic acid. Temperature. Mental state. Changes in the electrocardiogram. Prognosis If the vital organs are deprived of perfusion for more than just a short time, the prognosis is generally not good. Shock is still a medical emergency characterized by a high mortality rate. Early identification of patients who are likely to succumb to their illness is of utmost importance. Epidemiology Blood loss Trauma remains a leading cause of death worldwide with approximately half of these attributed to hemorrhage. In the United States in 2001, trauma was the third leading cause of death overall, and the leading cause of death in those aged 1 to 44 years. While trauma spans all demographics, it disproportionately affects the young with 40% of injuries occurring in ages 20 to 39 years by one country's account. Of this 40%, the greatest incidence was in the 20 to 24-year-old range. The preponderance of hemorrhagic shock cases resulting from trauma is high. During one year, one trauma center reported 62.2% of massive transfusions occur in the setting of trauma. The remaining cases are divided among cardiovascular surgery, critical care, cardiology, obstetrics, and general surgery, with trauma utilizing over 75% of the blood products. As patients age, physiological reserves decrease the likelihood of anticoagulant use increases and the number of comorbidities increases. Due to this, elderly patients are less likely to handle the physiological stresses of hemorrhagic shock and may decompensate more quickly. Fluid loss While the incidence of hypovolemic shock from extracellular fluid loss is difficult to quantify, it is known that hemorrhagic shock is most commonly due to trauma. In one study, 62.2% of massive transfusions at a level 1 trauma center were due to traumatic injury. In this study, 75% of the blood products used were related to traumatic injury. Elderly patients are more likely to experience hypovolemic shock due to fluid losses as they have less physiologic reserve. Hypovolemia secondary to diarrhea and/or dehydration is thought to be predominant in low-income countries.	Biology and health sciences	Cardiovascular disease	Health
960460	https://en.wikipedia.org/wiki/Tennis%20elbow	Tennis elbow	Tennis elbow, also known as lateral epicondylitis is an enthesopathy (attachment point disease) of the origin of the extensor carpi radialis brevis on the lateral epicondyle. It causes pain and tenderness over the bony part of the lateral epicondyle. Symptoms range from mild tenderness to severe, persistent pain. The pain may also extend into the back of the forearm. It usually has a gradual onset, but it can seem sudden and be misinterpreted as an injury. Tennis elbow is often idiopathic. Its cause and pathogenesis are unknown. It likely involves tendinosis, a degeneration of the local tendon. It is thought this condition is caused by excessive use of the muscles of the back of the forearm, but this is not supported by evidence. It may be associated with work or sports, classically racquet sports (including paddle sports), but most people with the condition are not exposed to these activities. The diagnosis is based on the symptoms and examination. Medical imaging is not very useful. Untreated enthesopathy usually resolves in 1–2 years. Treating the symptoms and pain involves medications such as NSAIDS or acetaminophen, a wrist brace, or a strap over the upper forearm. The role of corticosteroid injections as a form of treatment is still debated. Recent studies suggests that corticosteroid injections may delay symptom resolution. Signs and symptoms Patients typically feel pain or burning around the outer part of the elbow (lateral epicondyle of the humerus), which can move down the forearm and sometimes up to the upper arm. The pain is worsened by activities that involve wrist extension, such as gripping objects. Pain intensity varies from mild to severe and can be intermittent or constant, significantly impacting daily life. Patients also commonly report grip weakness and difficulty lifting. Terminology The term "tennis elbow" is widely used (although informal), but the condition affects non-tennis players. More recently, with the explosive growth of pickleball, the term "pickleball elbow" is frequently used. Historically, the medical term "lateral epicondylitis" was most commonly used for the condition, but "itis" implies inflammation and the condition is not inflammatory. It is also referred to as enthesopathy of the extensor carpi radialis origin. Since histological findings reveal noninflammatory tissue, the terms "lateral elbow tendinopathy" and "tendinosis" are suggested. In 2019, a group of international experts suggested that "lateral elbow tendinopathy" was the most appropriate terminology. But a disease of an attachment point (or enthesia) is most accurately referred to as an "enthesopathy." Causes The exact cause of lateral epicondylitis remains unclear. However, it is often linked to repetitive microtrauma resulting from excessive gripping, wrist extension, radial deviation, and/or forearm supination. Traditionally, people have speculated that tennis elbow is a type of repetitive strain injury resulting from tendon overuse and failed healing of the tendon, but there is no evidence of injury or repair, and misinterpretation of painful activities as a source of damage is common. Pathophysiology The extensor carpi radialis brevis is the most commonly affected muscle in lateral epicondylitis (LE), along with other extensor carpal muscles. Due to its unique origin, the ECRB tendon is prone to abrasion during elbow movements, leading to repetitive microtrauma. Lateral epicondylitis was initially considered an inflammatory process, however there is no evidence of inflammation or repair. Therefore, the disorder is more appropriately referred to as tendinosis or tendinopopathy. Tendinosis, a degenerative condition with fibroblasts, abnormal collagen, and increased blood vessels. Repetitive stress causes microtears, scar tissue formation, and biomechanical changes, worsening symptoms over time. Recently, successful results of a prospective therapeutic study of tennis elbow were published. It was observed that tennis elbow symptoms were most painful after awakening. It was hypothesized that a very common sleep position was interfering with healing and causing pain. The study evaluated if changing this position would avoid pressure on the lateral elbow while asleep. Patients who changed this sleep position reported successful resolution of symptoms, whereas those who were unable to change continued to have pain. The conclusion reached is that the pathophysiology of tennis elbow is due to an initial microscopic tear from a sprain/strain. This initial injury is aggravated at night by pressure on the sprain which delays healing. In other words, tennis elbow is neither a tendonitis nor a tendinosis, but more like a pressure sore. If the pressure is removed the initial injury goes on to heal. The importance of this finding is that other conditions characterized by nocturnal or early morning symptoms may also be worsened by a “pathological sleep position.” We know this applies to carpal and cubital tunnel syndrome, plantar fasciitis, shoulder/neck pain and Gerd. Clinical evaluation Physical examination Diagnosis is based on symptoms and clinical signs that are discrete and characteristic. For example, the extension of the elbow and flexion of the wrist causes outer elbow pain. The physical examination usually reveals marked tenderness at the origin of the extensor carpi radialis brevis muscle from the lateral epicondyle (extensor carpi radialis brevis origin). Pain may worsen with resisted wrist extension, middle finger extension, and forearm supination with an extended elbow, although normal elbow movement is often maintained, even in severe cases. Cozen's test Cozen's test is a physical examination performed to evaluate for tennis elbow involving pain with resisted wrist extension. The test is said to be positive if a resisted wrist extension triggers pain to the lateral aspect of the elbow owing to stress placed upon the tendon of the extensor carpi radialis brevis muscle. The test is performed with extended elbow. NOTE: With elbow flexed the extensor carpi radialis longus is in a shortened position as its origin is the lateral supracondylar ridge of the humerus. To rule out the ECRB (extensor carpi radialis brevis), repeat the test with the elbow in full extension. Medical imaging Medical imaging is not necessary or helpful. Radiographs (X-rays) may demonstrate calcification where the extensor muscles attach to the lateral epicondyle. Medical ultrasonography and magnetic resonance imaging (MRI) can demonstrate the pathology, but are not helpful for diagnosis and do not influence treatment. Longitudinal sonogram of the lateral elbow displays thickening and heterogeneity of the common extensor tendon that is consistent with tendinosis, as the ultrasound reveals calcifications, intrasubstance tears, and marked irregularity of the lateral epicondyle. Although the term “epicondylitis” is frequently used to describe this disorder, most histopathologic findings of studies have displayed no evidence of an acute, or a chronic inflammatory process. Histologic studies have demonstrated that this condition is the result of tendon degeneration, which replaces normal tissue with a disorganized arrangement of collagen. Colour Doppler ultrasound reveals structural tendon changes, with vascularity and hypo-echoic areas that correspond to the areas of pain in the extensor origin. Table of Clinical classification of lateral epicondylitis phases. Prevention Activity modification is the best way to prevent the occurrence of lateral epicondylitis. Prevention can include avoiding extreme end range motions in extension and flexion, limit repetitive hand and wrist motions, and modification of heavy lifting with extended arms. Lifestyle factors such as smoking, alcohol drinking, and dietary habits are known to influence the prognosis of various medical conditions. Smokers showed a higher chance of developing lateral epicondylitis compared to non-smokers. Current research indicates that alcohol intake is not significantly associated with lateral epicondylitis. Treatment Non-Operative Treatment Non operative treatment resolves 90% of symptomatic lateral epicondylitis. Nonoperative care usually includes activity modification, physical therapy, non-steroidal anti-inflammatory medications, bracing, extracorporeal shock-wave therapy, and acupuncture. Modifying activity and avoiding overuse are key to treatment. Lifting with the palm up and avoiding palm-down movements can shift strain from the lateral to the medial epicondyle, easing pain. Patients should also improve lifestyle habits and avoid triggering activities. Following the RICE method (rest, ice, compression, elevation) can help relieve pain initially. Exercises Stretching and isometric strengthening are the most common recommended exercises. The muscle is stretched with the elbow straight and the wrist passively flexed. Isometric strengthening can be done by pushing the top of the hand up against the undersurface of a table and holding the wrist straight. Orthotic devices Orthosis is a device externally used on the limb to improve the function or reduce the pain. Orthotics may be useful in tennis elbow; however, long-term effects are unknown. There are two main types of orthoses prescribed for this problem: counterforce elbow orthoses and wrist extension orthoses. Counterforce orthosis has a circumferential structure surrounding the arm. This orthosis usually has a strap which applies a binding force over the origin of the wrist extensors. The applied force by orthosis reduces the elongation within the musculotendinous fibers. Wrist extensor orthosis maintains the wrist in the slight extension. Speculative treatments Other approaches that are not experimentally tested include eccentric exercise using a rubber bar, joint manipulation directed at the elbow and wrist, spinal manipulation directed at the cervical and thoracic spinal regions, low level laser therapy, and extracorporeal shockwave therapy. Medication Recent studies demonstrate that topical nonsteroidal anti-inflammatory medications are effective within four weeks for lateral epicondylitis. Evidence for oral NSAIDs is mixed. Research indicates that corticosteroid injections improved outcomes more effectively than NSAIDs within four weeks but offered no long-term benefits at 12 months. Other studies suggest that, while helpful for short-term pain relief, corticosteroid injections are less effective than watchful waiting or physical therapy after one year. Repeated injections can also lead to tendon rupture and muscle atrophy. Thus, clinicians should be cautious with corticosteroid use for lateral epicondylitis due to limited long-term effectiveness and possible adverse effects. Alternative Treatments While many alternative treatments, such as shockwave, laser, low-frequency electrical nerve stimulation, ultrasound, and pulsed magnetic wave therapies, have been used, none have been proven effective. Current evidence is inconclusive on the effectiveness of acupuncture was effective for lateral epicondylitits. Platelet-Rich Plasma (PRP) Injections Platelet-Rich Plasma (PRP) has emerged as a potential treatment for lateral epicondylitis. PRP is derived from the patient's own blood and contains concentrated platelets, which are rich in growth factors. These growth factors are believed to initiate and accelerate tissue repair and regeneration support healing of the tendons and connective tissue and promote the growth of new blood vessels, aiding the recovery process. The PRP procedure for lateral epicondylitis involves extracting a small amount of the patient's blood, separating the plasma through centrifugation, and re-injecting it directly into the lateral epicondyle. While good outcomes have been reported with PRP for lateral epicondylitis, the overall literature is still unclear on its effectiveness. Additionally, variations in PRP preparation methods and injection techniques across different commercial systems add further complexity to assessing its effectiveness. Overall, current research on PRP as a treatment for lateral epicondylitis is promising. However, more studies are needed to provide clear evidence of its effectiveness. Surgery Most patients with lateral epicondylitis (tennis elbow) improve with conservative treatments and do not need surgery. However, if symptoms persist despite prolonged conservative therapy, surgical options should be reconsidered. Several surgical procedures are available for lateral epicondylitis, most involving the removal of damaged tissue from the ECRB and scraping of the lateral epicondyle. This procedure can be done through open, percutaneous, or arthroscopic methods. Percutaneous Surgery Percutaneous surgical approach is mainly used for releasing the common extensor tendon origin at the lateral epicondyle. This technique has been demonstrated to be safe, reliable, and cost-effective Good midterm outcomes in pain relief have been widely reported with a percutaneous surgical approach. However there is some limited evidence reported that arthroscopic and open techniques achieved a better prognosis than the percutaneous surgical approach for the treatment of lateral epicondylitis. In recent years, a new technique termed as ultrasound-guided percutaneous tenotomy has been reported as a safe and effective for the treatment of lateral epicondylitis, with improvements in symptoms, function, and ultrasound imaging at 1-year follow-up. Arthroscopic Surgery Arthroscopic surgery is a minimally invasive option for treating lateral epicondylitis. This technique fully visualizes the elbow joint, and leads to a quicker return to work. In the past, studies have shown good long term effects and fewer complications with arthroscopic surgery compared to open or percutaneous approaches. However, the literature is currently mixed with some recent reviews suggest no significant differences among open, arthroscopic, and percutaneous methods regarding recovery time, complication rates, or patient satisfaction. While others state that arthroscopic surgery may allow for a quicker return to work, suggesting a potential advantage in the early postoperative period. While results are generally positive, arthroscopic surgery carries risks of injury to the radial nerve and lateral ulnar collateral ligament. Epidemiology Tennis Elbow is a commonly seen condition and has been reported to affect 1% to 3% of adults each year. The incidence of lateral elbow tendinosis has declined, which could be due to shifts in diagnostic practices or an actual drop in cases. Understanding the typical disease progression can help patients and providers choose the best treatment approach. Symptoms of lateral epicondylitis Symptoms suggestive of lateral epicondylitis are present in about 1% of the adult population and are most common between ages 40 and 60. The prevalence varies somewhat between studies, likely as a result of varied diagnostic criteria and limited reliability between different observers. The data regarding symptoms of lateral epicondylitis in relation to occupations and sports are inconsistent and inconclusive. The shortcomings of the evidence that addresses the relationship between symptoms and occupation/sport include: variation in diagnostic criteria, limited reliability of diagnosis, confounding association of psychosocial factors, selection bias due to a high non-response rate, and the fact that exposures are usually by subjective patient reports and symptomatic patients might receive greater exposure. History German physician F. Runge is usually credited for the first description of the condition, calling it "writer's cramp" (Schreibekrampf) in 1873. Later, it was called as "washer women's elbow". British surgeon Henry Morris published an article in The Lancet describing "lawn tennis arm" in 1883. The popular term "tennis elbow" first appeared the same year in a paper by H. P. Major, described as "lawn-tennis elbow".	Biology and health sciences	Types	Health
960596	https://en.wikipedia.org/wiki/Messier%2022	Messier 22	Messier 22 or M22, also known as NGC 6656 or the Great Sagittarius Cluster, is an elliptical globular cluster of stars in the constellation Sagittarius, near the Galactic bulge region. It is one of the brightest globulars visible in the night sky. The brightest stars are 11th magnitude, with hundreds of stars bright enough to resolve with an 8" telescope. It is just south of the sun's position in mid-December, and northeast of Lambda Sagittarii (Kaus Borealis), the northernmost star of the "Teapot" asterism. M22 was one of the first globulars to be discovered, in 1665 by Abraham Ihle and it was included in Charles Messier's catalog of comet-like objects in 1764. It was one of the first globular clusters to be carefully studied – first by Harlow Shapley in 1930. He placed within it roughly 70,000 stars and found it had a dense core. Then Halton Arp and William G. Melbourne continued studies in 1959. Due to the large color spread of its red giant branch (RGB) sequence, akin to that in Omega Centauri, it became the object of intense scrutiny starting in 1977 with James E. Hesser et al. M22 is one of the nearer globular clusters to Earth – at about 10,600 light-years away. It spans 32′ on the sky which means its diameter (width across) is 99 ± 9 light-years, given its estimated distance. 32 variable stars have been recorded in M22. It is in front of part of the galactic bulge and is therefore useful for its microlensing effect on those background stars. Despite its relative proximity to us, this metal-poor cluster's light is limited by dust extinction, giving it an apparent magnitude of 5.5; even so, it is the brightest globular cluster visible from mid-northern latitudes (such as Japan, Korea, Europe and most of North America). From those latitudes due to its declination of nearly 24° south of the (celestial) equator, its daily path is low in the southern sky. It thus appears less impressive to people in the temperate northern hemisphere than counterparts fairly near in angle (best viewed in the Summer night sky) such as M13 and M5. M22 is one of only four globulars of our galaxy known to contain a planetary nebula (an expanding, glowing gas swell from a massive star, often a red giant). It was an object first noted of interest using the IRAS satellite by Fred Gillett and his associates in 1986, as a pointlike light source and its nature was found in 1989 by Gillett et al. The planetary nebula's central star is a blue star. The nebula, designated GJJC1, is likely about only 6,000 years old. Two black holes of between 10 and 20 solar masses () each were unearthed with the Very Large Array radio telescope in New Mexico and corroborated by the Chandra X-ray telescope in 2012. These imply that gravitational ejection of black holes from clusters is not as efficient as was previously thought, and leads to estimates of a total 5 to 100 black holes within M22. Interactions between stars and black holes could explain the unusually large core of the cluster. Gallery	Physical sciences	Notable star clusters	Astronomy
960789	https://en.wikipedia.org/wiki/Nitrous%20oxide%20%28medication%29	Nitrous oxide (medication)	Nitrous oxide, as medical gas supply, is an inhaled gas used as pain medication, and is typically administered with 50% oxygen mix. It is often used together with other medications for anesthesia. Common uses include during childbirth, following trauma, and as part of end-of-life care. Onset of effect is typically within half a minute, and the effect lasts for about a minute. Nitrous oxide was discovered between 1772 and 1793 and used for anesthesia in 1844. It is on the World Health Organization's List of Essential Medicines. It often comes as a 50/50 mixture with oxygen. Devices with a demand valve are available for self-administration. The setup and maintenance is relatively expensive for developing countries. There are few side effects, other than vomiting, with short-term use. With long-term use anemia or numbness may occur. It should always be given with at least 21% oxygen. It is not recommended in people with a bowel obstruction or pneumothorax. Use in the early part of pregnancy is not recommended. It is possible to continue breastfeeding following use. History Pure N2O was first used as a medical analgesic in December 1844, when Horace Wells made the first 12–15 dental operations with the gas in Hartford. Its debut as a generally accepted method, however, came in 1863, when Gardner Quincy Colton introduced it more broadly at all the Colton Dental Association clinics, that he founded in New Haven and New York City. The first devices used in dentistry to administer the gas consisted of a simple breathing bag made of rubber cloth. Breathing the pure gas often caused hypoxia (oxygen insufficiency) and sometimes death by asphyxiation. Eventually practitioners became aware of the need to provide at least 21% oxygen content in the gas (the same percentage as in air). In 1911, the anaesthetist Arthur Ernest Guedel first described the use of self-administration of a nitrous oxide and oxygen mix. It was not until 1961 that the first paper was published by Michael Tunstall and others, describing the administration of a pre-mixed 50:50 nitrous oxide and oxygen mix, which led to the commercialisation of the product. In 1970, Peter Baskett recognised that pre-mixed nitrous oxide and oxygen mix could have an important part to play in the provision of pre-hospital pain relief management, provided by ambulance personnel. Baskett contacted the Chief Ambulance Officer for the Gloucestershire Ambulance Brigade, Alan Withnell, to suggest this idea. This gained traction when Baskett negotiated with the British Oxygen Company, the availability of pre-mixed nitrous oxide and oxygen mix apparatus for training. Regular training sessions began at Frenchay Hospital (Bristol) and a pilot study was run in Gloucestershire (in which ambulances were crewed by a driver and one of the new highly trained ambulance men), the results of this trial were published in 1970. Today the nitrous oxide is administered in hospitals by a relative analgesia machine, which includes several improvements such as flowmeters and constant-flow regulators, an anaesthetic vaporiser, a medical ventilator, and a scavenger system, and delivers a precisely dosed and breath-actuated flow of nitrous oxide mixed with oxygen. The machine used in dentistry is much simpler, and is meant to be used by the patient in a fully conscious state. The gas is delivered through a demand-valve inhaler over the nose, which will only release gas when the patient inhales through it. Medical uses Nitrous oxide (N2O) is itself active (does not require any changes in the body to become active), and so has an onset in roughly the lung–brain circulation time with peak action 30 seconds after the start of administration. It is removed from the body unchanged via the lungs, and does not accumulate under normal conditions, explaining the rapid offset of around 60 seconds. It is effective in managing pain during labor and delivery. Nitrous oxide has been shown to be an effective and safe treatment for alcohol withdrawal. Nitrous oxide is more soluble than oxygen and nitrogen, so will tend to diffuse into any air spaces within the body. This makes it dangerous to use in patients with pneumothorax or those who have recently been scuba diving, and there are cautions over its use with any bowel obstruction. Its analgesic effect is strong (equivalent to 15 mg of subcutaneous route morphine) and characterised by rapid onset and offset, i.e. it is very fast-acting and wears off very quickly. When used in combination with other anesthetics gases, nitrous oxide causes a dose dependent increased respiratory rate and decreased tidal volumes, the net effect is a lower minute ventilation. Like volatile anesthetics, it increases cerebral blood flow and intracranial pressure. However, contrary to volatile anesthetics, it leads to an increase in cerebral metabolic rate of oxygen. Contraindications N2O should not be used in patients with bowel obstruction, pneumothorax, or middle ear or sinus disease, or who have had a recent intraocular injection of gas and should also not be used on any patient who has been scuba diving within the preceding 24 hours or in violently disturbed psychiatric patients. There are also clinical cautions in place for the first two trimesters of pregnancy and in patients with decreased levels of consciousness. Composition The gas is a mixture of half nitrous oxide (N2O) and half oxygen (O2). The ability to combine N2O and oxygen at high pressure while remaining in the gaseous form is caused by the Poynting effect (after John Henry Poynting, an English physicist). The Poynting effect involves the dissolution of gaseous O2 when bubbled through liquid N2O, with vaporisation of the liquid to form a gaseous O2/N2O mixture. Inhalation of pure N2O over a continued period would deprive the patient of oxygen, but the 50% oxygen content prevents this from occurring. The two gases will separate at low temperatures (<4 °C), which would permit administration of hypoxic mixtures. Therefore, it is not given from a cold cylinder without being shaken (usually by cylinder inversion) to remix the gases. Administration The gas is self-administered through a demand valve, using a mouthpiece, bite block or face mask. Self-administration of Entonox is safe because if enough is inhaled to start to induce anaesthesia, the patient becomes unable to hold the valve, and so will drop it and soon exhale the residual gas. This means that unlike other anaesthetic gases, it does not require the presence of an anaesthetist for administration. The 50% oxygen in Entonox ensures the person will have sufficient oxygen in their alveoli and conducting airways for a short period of apnea to be safe. Mechanism of action The pharmacological mechanism of action of in medicine is not fully known. However, it has been shown to directly modulate a broad range of ligand-gated ion channels, and this likely plays a major role in many of its effects. It moderately blocks NMDAR and β-subunit-containing nACh channels, weakly inhibits AMPA, kainate, GABA and 5-HT receptors, and slightly potentiates GABA and glycine receptors. It also has been shown to activate two-pore-domain channels. While affects quite a few ion channels, its anesthetic, hallucinogenic and euphoriant effects are likely caused predominantly, or fully, via inhibition of NMDA receptor-mediated currents. In addition to its effects on ion channels, may act to imitate nitric oxide (NO) in the central nervous system, and this may be related to its analgesic and anxiolytic properties. Nitrous oxide is 30 to 40 times more soluble than nitrogen. Society and culture Nitronox was a registered trademark of the BOC Group between 1966 and 1999, and was reregistered by Hs Tm Inc since 2005 It is also colloquially known as "gas and air" in the United Kingdom. Research Investigational trials show potential for antidepressant applications of N2O, especially for treatment-resistant forms of depression, and it is rapid-acting. In a phase 2 clinical trial, a treatment with 25% nitrous oxide had comparable efficacy to 50% nitrous oxide but was associated with significantly fewer adverse effects.	Biology and health sciences	Anesthetics	Health
961349	https://en.wikipedia.org/wiki/Hot%20Jupiter	Hot Jupiter	Hot Jupiters (sometimes called hot Saturns) are a class of gas giant exoplanets that are inferred to be physically similar to Jupiter but that have very short orbital periods (). The close proximity to their stars and high surface-atmosphere temperatures resulted in their informal name "hot Jupiters". Hot Jupiters are the easiest extrasolar planets to detect via the radial-velocity method, because the oscillations they induce in their parent stars' motion are relatively large and rapid compared to those of other known types of planets. One of the best-known hot Jupiters is . Discovered in 1995, it was the first extrasolar planet found orbiting a Sun-like star. has an orbital period of about four days. General characteristics Though there is diversity among hot Jupiters, they do share some common properties. Their defining characteristics are their large masses and short orbital periods, spanning 0.36–11.8 Jupiter masses and 1.3–111 Earth days. The mass cannot be greater than approximately 13.6 Jupiter masses because then the pressure and temperature inside the planet would be high enough to cause deuterium fusion, and the planet would be a brown dwarf. Most have nearly circular orbits (low eccentricities). It is thought that their orbits are circularized by perturbations from nearby stars or tidal forces. Whether they remain in these circular orbits for long periods of time or collide with their host stars depends on the coupling of their orbital and physical evolution, which are related through the dissipation of energy and tidal deformation. Many have unusually low densities. The lowest one measured thus far is that of TrES-4b at 0.222 g/cm3. The large radii of hot Jupiters are not yet fully understood but it is thought that the expanded envelopes can be attributed to high stellar irradiation, high atmospheric opacities, possible internal energy sources, and orbits close enough to their stars for the outer layers of the planets to exceed their Roche limit and be pulled further outward. Usually they are tidally locked, with one side always facing its host star. They are likely to have extreme and exotic atmospheres due to their short periods, relatively long days, and tidal locking. Atmospheric dynamics models predict strong vertical stratification with intense winds and super-rotating equatorial jets driven by radiative forcing and the transfer of heat and momentum. Recent models also predict a variety of storms (vortices) that can mix their atmospheres and transport hot and cold regions of gas. The day-night temperature difference at the photosphere is predicted to be substantial, approximately for a model based on HD 209458 b. They appear to be more common around F- and G-type stars and less so around K-type stars. Hot Jupiters around red dwarfs are very rare. Generalizations about the distribution of these planets must take into account the various observational biases, but in general their prevalence decreases exponentially as a function of the absolute stellar magnitude. Formation and evolution There are three schools of thought regarding the possible origin of hot Jupiters. One possibility is that they were formed in situ at the distances at which they are currently observed. Another possibility is that they were formed at a distance but later migrated inward. Such a shift in position might occur due to interactions with gas and dust during the solar nebula phase. It might also occur as a result of a close encounter with another large object destabilizing a Jupiter's orbit. Migration In the migration hypothesis, a hot Jupiter forms beyond the frost line, from rock, ice, and gases via the core accretion method of planetary formation. The planet then migrates inwards to the star where it eventually forms a stable orbit. The planet may have migrated inward smoothly via type II orbital migration. Or it may have migrated more suddenly due to gravitational scattering onto eccentric orbits during an encounter with another massive planet, followed by the circularization and shrinking of the orbits due to tidal interactions with the star. A hot Jupiter's orbit could also have been altered via the Kozai mechanism, causing an exchange of inclination for eccentricity resulting in a high eccentricity low perihelion orbit, in combination with tidal friction. This requires a massive body—another planet or a stellar companion—on a more distant and inclined orbit; approximately 50% of hot Jupiters have distant Jupiter-mass or larger companions, which can leave the hot Jupiter with an orbit inclined relative to the star's rotation. The type II migration happens during the solar nebula phase, i.e. when gas is still present. Energetic stellar photons and strong stellar winds at this time remove most of the remaining nebula. Migration via the other mechanism can happen after the loss of the gas disk. In situ Instead of being gas giants that migrated inward, in an alternate hypothesis the cores of the hot Jupiters began as more common super-Earths which accreted their gas envelopes at their current locations, becoming gas giants in situ. The super-Earths providing the cores in this hypothesis could have formed either in situ or at greater distances and have undergone migration before acquiring their gas envelopes. Since super-Earths are often found with companions, the hot Jupiters formed in situ could also be expected to have companions. The increase of the mass of the locally growing hot Jupiter has a number of possible effects on neighboring planets. If the hot Jupiter maintains an eccentricity greater than 0.01, sweeping secular resonances can increase the eccentricity of a companion planet, causing it to collide with the hot Jupiter. The core of the hot Jupiter in this case would be unusually large. If the hot Jupiter's eccentricity remains small the sweeping secular resonances could also tilt the orbit of the companion. Traditionally, the in situ mode of conglomeration has been disfavored because the assembly of massive cores, which is necessary for the formation of hot Jupiters, requires surface densities of solids ≈ 104 g/cm2, or larger. Recent surveys, however, have found that the inner regions of planetary systems are frequently occupied by super-Earth type planets. If these super-Earths formed at greater distances and migrated closer, the formation of in situ hot Jupiters is not entirely in situ. Atmospheric loss If the atmosphere of a hot Jupiter is stripped away via hydrodynamic escape, its core may become a chthonian planet. The amount of gas removed from the outermost layers depends on the planet's size, the gases forming the envelope, the orbital distance from the star, and the star's luminosity. In a typical system, a gas giant orbiting at 0.02 AU around its parent star loses 5–7% of its mass during its lifetime, but orbiting closer than 0.015 AU can mean evaporation of a substantially larger fraction of the planet's mass. No such objects have been found yet and they are still hypothetical. Terrestrial planets in systems with hot Jupiters Simulations have shown that the migration of a Jupiter-sized planet through the inner protoplanetary disk (the region between 5 and 0.1 AU from the star) is not as destructive as expected. More than 60% of the solid disk materials in that region are scattered outward, including planetesimals and protoplanets, allowing the planet-forming disk to reform in the gas giant's wake. In the simulation, planets up to two Earth masses were able to form in the habitable zone after the hot Jupiter passed through and its orbit stabilized at 0.1 AU. Due to the mixing of inner-planetary-system material with outer-planetary-system material from beyond the frost line, simulations indicated that the terrestrial planets that formed after a hot Jupiter's passage would be particularly water-rich. According to a 2011 study, hot Jupiters may become disrupted planets while migrating inwards; this could explain an abundance of "hot" Earth-sized to Neptune-sized planets within 0.2 AU of their host star. One example of these sorts of systems is that of WASP-47. There are three inner planets and an outer gas giant in the habitable zone. The innermost planet, WASP-47e, is a large terrestrial planet of 6.83 Earth masses and 1.8 Earth radii; the hot Jupiter, b, is little heavier than Jupiter, but about 12.63 Earth radii; a final hot Neptune, c, is 15.2 Earth masses and 3.6 Earth radii. A similar orbital architecture is also exhibited by the Kepler-30 system. Misaligned orbits Several hot Jupiters, such as HD 80606 b, have orbits that are misaligned with their host stars, including several with retrograde orbits such as HAT-P-14b. This misalignment may be related to the heat of the photosphere the hot Jupiter is orbiting. There are several proposed hypotheses as to why this might occur. One such hypothesis involves tidal dissipation and suggests there is a single mechanism for producing hot Jupiters and this mechanism yields a range of obliquities. Cooler stars with higher tidal dissipation damps the obliquity (explaining why hot Jupiters orbiting cooler stars are well aligned) while hotter stars do not damp the obliquity (explaining the observed misalignment). Another hypothesis is that the host star sometimes changes rotation early in its evolution, rather than the orbit changing. Yet another hypothesis is that hot Jupiters tend to form in dense clusters, where perturbations are more common and gravitational capture of planets by neighboring stars is possible. Ultra-hot Jupiters Ultra-hot Jupiters are hot Jupiters with a dayside temperature greater than . In such dayside atmospheres, most molecules dissociate into their constituent atoms and circulate to the nightside where they recombine into molecules again. One example is TOI-1431b, announced by the University of Southern Queensland in April 2021, which has an orbital period of just two and a half days. Its dayside temperature is , making it hotter than 40% of stars in our galaxy. The nightside temperature is . Ultra-short period planets Ultra-short period planets (USP) are a class of planets with orbital periods below one day and occur only around stars of less than about 1.25 solar masses. Confirmed transiting hot Jupiters that have orbital periods of less than one day include WASP-18b, Banksia, Astrolábos and WASP-103b. Puffy planets Gas giants with a large radius and very low density are sometimes called "puffy planets" or "hot Saturns", due to their density being similar to Saturn's. Puffy planets orbit close to their stars so that the intense heat from the star combined with internal heating within the planet will help inflate the atmosphere. Six large-radius low-density planets have been detected by the transit method. In order of discovery they are: HAT-P-1b, CoRoT-1b, TrES-4b, WASP-12b, WASP-17b, and Kepler-7b. Some hot Jupiters detected by the radial-velocity method may be puffy planets. Most of these planets are around or below Jupiter mass as more massive planets have stronger gravity keeping them at roughly Jupiter's size. Indeed, hot Jupiters with masses below Jupiter, and temperatures above 1800 Kelvin, are so inflated and puffed out that they are all on unstable evolutionary paths which eventually lead to Roche-Lobe overflow and the evaporation and loss of the planet's atmosphere. Even when taking surface heating from the star into account, many transiting hot Jupiters have a larger radius than expected. This could be caused by the interaction between atmospheric winds and the planet's magnetosphere creating an electric current through the planet that heats it up, causing it to expand. The hotter the planet, the greater the atmospheric ionization, and thus the greater the magnitude of the interaction and the larger the electric current, leading to more heating and expansion of the planet. This theory matches the observation that planetary temperature is correlated with inflated planetary radii. Moons Theoretical research suggests that hot Jupiters are unlikely to have moons, due to both a small Hill sphere and the tidal forces of the stars they orbit, which would destabilize any satellite's orbit, the latter process being stronger for larger moons. This means that for most hot Jupiters, stable satellites would be small asteroid-sized bodies. Furthermore, the physical evolution of hot Jupiters can determine the final fate of their moons: stall them in semi-asymptotic semimajor axes, or eject them from the system where they may undergo other unknown processes. In spite of this, observations of WASP-12b suggest that it is orbited by at least one large exomoon. Hot Jupiters around red giants It has been proposed that gas giants orbiting red giants at distances similar to that of Jupiter could be hot Jupiters due to the intense irradiation they would receive from their stars. It is very likely that in the Solar System Jupiter will become a hot Jupiter after the transformation of the Sun into a red giant. The recent discovery of particularly low-density gas giants orbiting red giant stars supports this hypothesis. Hot Jupiters orbiting red giants would differ from those orbiting main-sequence stars in a number of ways, most notably the possibility of accreting material from the stellar winds of their stars and, assuming a fast rotation (not tidally locked to their stars), a much more evenly distributed heat with many narrow-banded jets. Their detection using the transit method would be much more difficult due to their tiny size compared to the stars they orbit, as well as the long time needed (months or even years) for one to transit their star as well as to be occulted by it. Star–planet interactions Theoretical research since 2000 suggested that "hot Jupiters" may cause increased flaring due to the interaction of the magnetic fields of the star and its orbiting exoplanet, or because of tidal forces between them. These effects are called "star–planet interactions" or SPIs. The HD 189733 system is the best-studied exoplanet system where this effect was thought to occur. In 2008, a team of astronomers first described how as the exoplanet orbiting HD 189733 A reaches a certain place in its orbit, it causes increased stellar flaring. In 2010, a different team found that every time they observe the exoplanet at a certain position in its orbit, they also detected X-ray flares. In 2019, astronomers analyzed data from Arecibo Observatory, MOST, and the Automated Photoelectric Telescope, in addition to historical observations of the star at radio, optical, ultraviolet, and X-ray wavelengths to examine these claims. They found that the previous claims were exaggerated and the host star failed to display many of the brightness and spectral characteristics associated with stellar flaring and solar active regions, including sunspots. Their statistical analysis also found that many stellar flares are seen regardless of the position of the exoplanet, therefore debunking the earlier claims. The magnetic fields of the host star and exoplanet do not interact, and this system is no longer believed to have a "star-planet interaction." Some researchers had also suggested that HD 189733 accretes, or pulls, material from its orbiting exoplanet at a rate similar to those found around young protostars in T Tauri star systems. Later analysis demonstrated that very little, if any, gas was accreted from the "hot Jupiter" companion.	Physical sciences	Planetary science	Astronomy
961611	https://en.wikipedia.org/wiki/Lip%20balm	Lip balm	Lip balm or lip salve is a wax-like substance applied to the lips to moisturize and relieve chapped or dry lips, angular cheilitis, stomatitis, or cold sores. Lip balm often contains beeswax or carnauba wax, camphor, cetyl alcohol, lanolin, paraffin, and petrolatum, among other ingredients. Some varieties contain dyes, flavor, fragrance, phenol, salicylic acid, and sunscreen. Overview The primary purpose of lip balm is to provide an occlusive layer on the lip surface to seal moisture in lips and protect them from external exposure. Dry air, cold temperatures, and wind all have a drying effect on skin by drawing moisture away from the body. Lips are particularly vulnerable because the skin is so thin, and thus they are often the first to present signs of dryness. Occlusive materials like waxes and petroleum jelly prevent moisture loss and maintain lip comfort while flavorings, colorants, sunscreens, and various medicaments can provide additional, specific benefits. Lip balms are produced from bee wax and natural candelilla and carnauba waxes. Lip balm can be applied by a finger to the lips, or in a lipstick-style tube from which it can be applied directly. In 2022, the global lip balm market was valued at US$732.76 mln. The market is predicted to grow at a rate of 9.28% within the next five years and is likely to reach US$1247.74 mln by 2027. Production Production for lip balms includes the following stages: Raw materials are checked for its quality (cosmetic products must comply with strict safety standards) The ingredients are dosed, melted, and mixed (this stage involves special equipment) This mixture is treated in a vacuum to remove bubbles The mixture is crystallized for about 48 hours The mixture is then remelted The mixture is cut into pieces which are shaped as required The lip balm is packaged into a casing History Early lip balms Since 40 BC, the Egyptians made treatment for lip care, which was made with a mixture of beeswax, olive oil, and animal fat. United States In the 1800s, Lydia Maria Child recommended earwax as a treatment for cracked lips in her highly-popular book, Child observed that, "Those who are troubled with cracked lips have found this earwax remedy successful when others have failed. It is one of those sorts of cures, which are very likely to be laughed at; but I know of its having produced very beneficial results." The invention of the lip balm was first formally invented in the 1880s by physician Charles Brown Fleet though its origins may be traced to earwax. Fleet later named his lip balm product "ChapStick". In 1872, chemist Robert Chesebrough discovered and sampled a new petroleum jelly, initially describing it as a "natural, waxy ingredient, rich in minerals from deep within the earth" which could be used as a solution for skin repair. He then distributed his product under the name "Wonder Jelly" before shortly changing it to "Vaseline". In the early 1880s, Charles Brown Fleet created ChapStick. However, due to the lack of sales, Fleet sold his formula and rights to ChapStick to John Morton in 1912 for $5, who saw the marketing potential in the brand. After making the purchase, Morton commissioned Frank Wright, Jr. to create a design for the logo of ChapStick for $15 in 1936. In 1972, ChapStick tubes concealing hidden microphones were used during the Watergate scandal. In 1937, Alfred Woelbing created Carmex to treat cold sores in Milwaukee, though the occurrence of World War 2 would slow the production and sales due to the lack of lanolin. In 1980, Carmex underwent a product change by converting its packaging into squeezable tubes. In 1973, Bonne Bell created the first flavored lip balm and marketed the company as Lip Smackers. The company would later collaborate on various different-flavored lip balms including Dr. Pepper in 1975, The Wrigley Company in 2004, and The Coca-Cola Company in 2006. Bonne Bell also collaborated with Disney to produce lip balms with various princess characters in 2010. In 1991, Burt Shavitz and Roxanne Quimby created their first beeswax based lip balm solution through their company, Burt's Bees. In 2020, it was reported that Burt's Bees had used 50 percent of recycled material to package various products and that 100 percent of the products were recyclable. In 2011, Evolution of Smooth (or commonly known as EOS) created a spherical-shaped lip balm as well as describing its 95% organic ingredients. Cannabis infused lip balms With the gradual legalization of cannabis in the United States, some companies have produced lip balms containing doses of THC or CBD oil. The lip balms were infused with a low dosage of THC in order to prevent the occurrence of any psychoactive or related effect. Notable brands Burt's Bees Blistex Carmex ChapStick Labello Lip Smacker Lypsyl EOS Vaseline Aquaphor Nivea Dependency Addictive ingredients Some physicians have suggested that certain types of lip balm can be addictive or contain ingredients that actually cause drying, the accuracy of which has been debated by many professionals. Lip balm manufacturers sometimes state in their FAQs that there is nothing addictive in their products or that all ingredients are listed and approved by the FDA. Snopes found the claim that there are substances in Carmex that are irritants necessitating reapplication, such as ground glass, to be false. However, some experts such as dermatologist Dr. Cynthia Bailey state that some ingredients in lip balm directly causes sensitive lip skin which may lead to addiction. Dermatology professor Marcia Driscoll also adds onto this argument by stating that aroma ingredients found in flavored or scented lip balms have the potential to irritate skin. Causes for Dependency According to a report, professor Brad Rohu states that it is natural for the lips to feel dry. The exposure to environments with cold, dry, or windy weather can directly cause the chapping of the lips as well as behaviors such as lip licking or mouth breathing. These factors may directly contribute to an increased amount of lip balm usage. According to dermatologist Amy Derick, those who have expressed dependencies on lip balm have developed a desire of how the lips feel after application. She also mentions that the variety of lip balm flavor may also directly cause lip balm dependency as a person may want to lick their lips to taste the flavor, which may consequentially remove the lip balm coating from the lips. This may also leave saliva on the lips which can dry up and make the lips feel even more dry than they initially were. Effects on lip barrier The human lips have an inadequate capability of holding moisture as well as an imperfect lip barrier function. The Journal of the American Academy of Dermatology performed a study in order to determine whether consistent use of lip balm would enhance the overall quality of the lips. The study used 32 female participants within the ages of 20 to 40 years and the participants had mild to moderate dried lips without any history of health-related complications. The participants underwent a procedure in which no lip treatment was provided on the first 3 days, then 2 weeks of consistent lip balm usage, and then a period of no treatment for 3 days. The study determined the quality of the lips based on the physical details and appearance throughout the study. The study showed a direct improvement of the physical details of the lips except for lip cracking during the second week of treatment and after the period of no treatment. The study also showed that hydration of the lips lasted for approximately 8 hours after usage and the lip balm improved the lip barrier function despite discontinued usage. The study concluded that lip balms assist the hydration of the lips which consequentially improves the lip barrier function and the quality. This study was completely funded by Burt's Bees, a lip balm company. Mineral oil In 2015, German consumer watchdog Stiftung Warentest analyzed cosmetics containing mineral oils. After developing a new detection method they found high concentrations of Mineral Oil Aromatic Hydrocarbons (MOAH) and even polyaromatics in products containing mineral oils with Vaseline products containing the most MOAH of all tested cosmetics (up to 9%). The European Food Safety Authority sees MOAH and polyaromatics as possibly carcinogenic. Based on the results, Stiftung Warentest warns not to use Vaseline or any product that is based on mineral oils for lip care. Lip balm market United States In 2019, a research report conducted by the Statista Research Department concluded that ChapStick was the leading lip balm brand in the United States with an approximate unit sale of 55.8 million. Carmex was the second leading brand with approximately 35.2 million units sold and Burt's Bees being the third leading brand with approximately 32.3 million units sold. Trends Beezin' Beezin' is a trend dating back to 2013 in which a person applies Burt's Bees brand lip balm onto the eyelids. The practice is done in order to feel a sensation of being high or drunk, and even to increase the desired effects of alcohol and other substances. In 2022, Beezin' became a viral trend on the social media platform TikTok. Some ingredients, including peppermint oil, are known to be eye irritants which can cause an unintentional inflammatory response which may require treatment and may also cause dermatitis on the eyelids.	Biology and health sciences	Hygiene products	Health
961680	https://en.wikipedia.org/wiki/Malayan%20tapir	Malayan tapir	The Malayan tapir (Tapirus indicus), also called Asian tapir, Asiatic tapir, oriental tapir, Indian tapir, piebald tapir, or black-and-white tapir, is the only living tapir species outside of the Americas. It is native to Southeast Asia from the Malay Peninsula to Sumatra. It has been listed as Endangered on the IUCN Red List since 2008, as the population is estimated to comprise fewer than 2,500 mature individuals. Taxonomy The scientific name Tapirus indicus was proposed by Anselme Gaëtan Desmarest in 1819 who referred to a tapir described by Pierre-Médard Diard. Tapirus indicus brevetianus was coined by a Dutch zoologist in 1926 who described a black Malayan tapir from Sumatra that had been sent to Rotterdam Zoo in the early 1920s. Phylogenetic analyses of 13 Malayan tapirs showed that the species is monophyletic. It was placed in the genus Acrocodia by Colin Groves and Peter Grubb in 2011. However, a comparison of mitochondrial DNA of 16 perissodactyl species revealed that the Malayan tapir forms a sister group together with the Tapirus species native to the Americas. It was the first Tapirus species that genetically diverged from the group, estimated about in the Late Oligocene. Description The Malayan tapir is easily identified by its markings, most notably the light-colored patch that extends from its shoulders to its hindquarters. Black hair covers its head, shoulders, and legs, while white hair covers its midsection, rear, and the tips of its ears; these white edges around the rims of the outer ear as is true of other tapirs. The disrupted coloration breaks up its outline, providing camouflage by making the animal difficult to recognize against the varied terrain and dense flora of its habitat; potential predators may mistake it for a large rock, rather than prey, when it is lying down to sleep. The Malayan tapir is the largest of the four extant tapir species and grows to between in length, not counting a stubby tail of only in length, and stands tall. It typically weighs between , although some adults can weigh up to . The females are usually larger than the males. Like other tapir species, it has a small, stubby tail and a long, flexible proboscis. It has four toes on each front foot and three toes on each back foot. The Malayan tapir has rather poor eyesight, but excellent hearing and sense of smell. The tapir's unique proboscis is supported by several evolutionary adaptations of its skull. It has a large sagittal crest, unusually positioned orbits, an unusually shaped cranium with elevated frontal bones, and a retracted nasal incision as well as retracted facial cartilage. This evolutionary process is believed to have caused the loss of some cartilages, facial muscles, and the bony wall of the tapir's nasal chamber. Vision Malayan tapirs have very poor eyesight, both on land and in water, instead relying heavily on their excellent senses of smell and hearing to navigate and forage. Their eyes are small and, like many herbivores, positioned on the sides of the face. They have brown irises, but the corneas are often covered in a blue haze; this corneal cloudiness is thought to be caused by repetitive exposure to light. This loss of transparency impacts the ability of the cornea to transmit and focus outside light as it enters the eye, impairing the animal's overall vision. As these tapirs are most active at night on top of having poor eyesight, this habit may make it harder for them to search for food and avoid predators. Color variation Two melanistic Malayan tapirs were observed in Jerangau Forest Reserve in Malaysia in 2000. A black Malayan tapir was also recorded in Tekai Tembeling Forest Reserve in Pahang state in 2016. Distribution and habitat The Malayan tapir lives throughout the tropical lowland rainforests of Southeast Asia, including Sumatra in Indonesia, Peninsular Malaysia, Myanmar, and Thailand. Pleistocene fossils were found in Java and other locations accompanied by herbivores more typical of grasslands, indicating that it evolved in more open habitats and retreated to closed forests in later times. It was found in Borneo until at least 8,000 years ago during the early Holocene in the Niah Caves of Sarawak, and some 19th century writers mentioned it as a contemporary species in Borneo, likely based on native accounts. It has been proposed to reintroduce the tapir to the island as a conservation measure. In the continent, the Malayan tapir was found in historical times as far north as China. Behaviour and ecology Malayan tapirs are primarily solitary, marking out large tracts of land as their territory, though these areas usually overlap with those of other individuals. Tapirs mark out their territories by spraying urine on plants, and they often follow distinct paths which they have bulldozed through the undergrowth. Exclusively herbivorous, the animal forages for the tender shoots and leaves of more than 115 species of plants, of which around 30 are particularly preferred, moving slowly through the forest and pausing often to eat and note the scents left behind by other tapirs in the area. The tapir can run quickly when threatened or frightened, and if forced to fight can defend itself with its strong jaws and sharp teeth. Malayan tapirs communicate with high-pitched squeaks and whistles. They usually prefer to live near water and often bathe and swim, and they are also able to climb steep slopes. Tapirs are mainly active at night, though they are not exclusively nocturnal; because they tend to eat soon after sunset or before sunrise, and they will often nap in the middle of the night, they are considered to be crepuscular animals. Life cycle The gestation period of the Malayan tapir is about 390–395 days, after which a single calf is born that weighs around . Malayan tapirs are the largest of the four tapir species at birth and tend to grow more quickly than their relatives. Young tapirs of all species have brown hair with white stripes and spots, a pattern that enables them to hide effectively in the dappled light of the forest. This baby coat fades into adult coloration between four and seven months after birth. Weaning occurs between six and eight months of age, at which time the babies are nearly full-grown, and the animals reach sexual maturity around age three. Breeding typically occurs in April, May or June, and females generally produce one calf every two years. Malayan tapirs can live up to 30 years, both in the wild and in captivity. Predators Because of its size, the Malayan tapir has few natural predators, and even reports of killings by tigers (Panthera tigris) are scarce. Malayan tapirs can defend themselves with their very powerful bite; in 1998, the bite of a captive female Malayan tapir severed off a zookeeper's left arm at the mid-bicep, likely because she stood between her and her offspring. Threats The main threats to the Malayan tapir are loss and destruction of habitat through deforestation. Large tracts of forests in Thailand and Malaysia have been converted for planting oil palms. Habitat fragmentation in peninsular Malaysia caused displacement of 142 Malayan tapirs between 2006 and 2010; some were rescued and relocated, while 15 of them were killed in vehicle collisions.	Biology and health sciences	Perissodactyla	Animals
961724	https://en.wikipedia.org/wiki/Choy%20sum	Choy sum	Choy sum (also spelled choi sum, choi sam in Cantonese; cai xin, caixin in Standard Mandarin) is a leafy vegetable commonly used in Chinese cuisine. It is a member of the genus Brassica of the mustard family, Brassicaceae (Brassica rapa var. parachinensis or Brassica chinensis var. parachinensis). Choy sum is a transliteration of the Cantonese name (), which can be literally translated as "heart of the vegetable". Choy sum is also called yu choy (you cai in Standard Mandarin; Chinese: 油菜). It is also known as Chinese flowering cabbage. Description Choy sum is a green leafy vegetable similar to gai lan, and can be characterized by the distinct yellow flowers which it bears. Each flower has four yellow, oval to round petals with six stamens on fleshy, erect stems which are in diameter and tall with light to dark green, and are oval (becomes acuminate shaped, or basal-shaped near the flowering stage) with slightly serrated margins leaves, which never forms compact heads like the cabbage. Fruits can develop out of cross-pollination or self-pollination, and are silique structured, that open at maturity through dehiscence or drying to bare open to brown or black seeds that are small and round in shape. A single pod can bear 4 to 46 seeds. The height of the plant varies greatly, ranging from depending on the growing conditions and the variety. Flowering usually appears when there are about 7 to 8 leaves on the plant or about tall. The bulk of the root system is found within a depth of and is confined to a radius of . The whole plant is overall an annual, herbaceous plant, rarely perennial, rarely growing into subshrubs. The whole plant consists of a simple or branched (when it is near the flowering stage), leafy structure. It grows best in soil with a minimum pH level of 5.6, maximum pH level of 7.5. Use Choy sum is highly valued as a vegetable in China and Japan. It is commonly consumed in soup, blanched, or stir-fried. Gallery	Biology and health sciences	Leafy vegetables	Plants
962035	https://en.wikipedia.org/wiki/Seismic%20refraction	Seismic refraction	Seismic refraction is a geophysical principle governed by Snell's Law of refraction. The seismic refraction method utilizes the refraction of seismic waves by rock or soil layers to characterize the subsurface geologic conditions and geologic structure. Seismic refraction is exploited in engineering geology, geotechnical engineering and exploration geophysics. Seismic refraction traverses (seismic lines) are performed using an array of seismographs or geophones and an energy source. The methods depend on the fact that seismic waves have differing velocities in different types of soil or rock. The waves are refracted when they cross the boundary between different types (or conditions) of soil or rock. The methods enable the general soil types and the approximate depth to strata boundaries, or to bedrock, to be determined. P-wave refraction P-wave refraction evaluates the compression wave generated by the seismic source located at a known distance from the array. The wave is generated by vertically striking a striker plate with a sledgehammer, shooting a seismic shotgun into the ground, or detonating an explosive charge in the ground. Since the compression wave is the fastest of the seismic waves, it is sometimes referred to as the primary wave and is usually more-readily identifiable within the seismic recording as compared to the other seismic waves. S-wave refraction S-wave refraction evaluates the shear wave generated by the seismic source located at a known distance from the array. The wave is generated by horizontally striking an object on the ground surface to induce the shear wave. Since the shear wave is the second fastest wave, it is sometimes referred to as the secondary wave. When compared to the compression wave, the shear wave is approximately one-half (but may vary significantly from this estimate) the velocity depending on the medium. Two horizontal layers ic0 - critical angle V0 - velocity of the first layer V1 - velocity of the second layer h0 - thickness of the first layer T01 - intercept Several horizontal layers Inversion methods The General Reciprocal method The Plus minus method Refraction inversion modeling (refraction tomography) Monte Carlo simulation Genetic algorithms Applications Seismic refraction has been successfully applied to tailings characterisation through P- and S-wave travel time tomographic inversions.	Physical sciences	Seismology	Earth science
962148	https://en.wikipedia.org/wiki/Beehive%20Cluster	Beehive Cluster	The Beehive Cluster (also known as Praesepe (Latin for "manger", "cot" or "crib"), M44, NGC 2632, or Cr 189), is an open cluster in the constellation Cancer. One of the nearest open clusters to Earth, it contains a larger population of stars than other nearby bright open clusters holding around 1,000 stars. Under dark skies, the Beehive Cluster looks like a small nebulous object to the naked eye, and has been known since ancient times. Classical astronomer Ptolemy described it as a "nebulous mass in the breast of Cancer". It was among the first objects that Galileo studied with his telescope. Age and proper motion coincide with those of the Hyades, suggesting they may share similar origins. Both clusters also contain red giants and white dwarfs, which represent later stages of stellar evolution, along with many main sequence stars. Distance to M44 is often cited to be between 160 and 187 parsecs (520–610 light years), but the revised Hipparcos parallaxes (2009) for Praesepe members and the latest infrared color-magnitude diagram favors an analogous distance of 182 pc. There are better age estimates of around 600 million years (compared to about 625 million years for the Hyades). The diameter of the bright inner cluster core is about 7.0 parsecs (23 light years). At 1.5° across, the cluster easily fits within the field of view of binoculars or low-powered small telescopes. Regulus, Castor, and Pollux are guide stars. History In 1609, Galileo first telescopically observed the Beehive and was able to resolve it into 40 stars. Charles Messier added it to his famous catalog in 1769 after precisely measuring its position in the sky. Along with the Orion Nebula and the Pleiades cluster, Messier's inclusion of the Beehive has been noted as curious, as most of Messier's objects were much fainter and more easily confused with comets. Another possibility is that Messier simply wanted to have a larger catalog than his scientific rival Lacaille, whose 1755 catalog contained 42 objects, and so he added some well-known bright objects to boost his list. Wilhelm Schur, as director of the Göttingen Observatory, drew a map of the cluster in 1894.Ancient Greeks and Romans saw this object as a manger from which two donkeys, the adjacent stars Asellus Borealis and Asellus Australis, are eating; these are the donkeys that Dionysos and Silenus rode into battle against the Titans. Hipparchus (c.130 BC) refers to the cluster as Nephelion ("Little Cloud") in his star catalog. Claudius Ptolemy's Almagest includes the Beehive Cluster as one of seven "nebulae" (four of which are real), describing it as "The Nebulous Mass in the Breast (of Cancer)". Aratus (c.260–270 BC) calls the cluster Achlus or "Little Mist" in his poem Phainomena. Johann Bayer showed the cluster as a nebulous star on his Uranometria atlas of 1603, and labeled it Epsilon. The letter is now applied specifically to the brightest star of the cluster Epsilon Cancri, of magnitude 6.29. This perceived nebulous object is in the Ghost (Gui Xiu), the 23rd lunar mansion of ancient Chinese astrology. Ancient Chinese skywatchers saw this as a ghost or demon riding in a carriage and likened its appearance to a "cloud of pollen blown from willow catkins". It was also known by the somewhat less romantic name of Jishi qi (積屍氣, also transliterated Tseih She Ke), the "Exhalation of Piled-up Corpses". It is also known simply as Jishi (積屍), "cumulative corpses". Morphology and composition Like many star clusters of all kinds, Praesepe has experienced mass segregation. This means that bright massive stars are concentrated in the cluster's core, while dimmer and less massive stars populate its halo (sometimes called the corona). The cluster's core radius is estimated at 3.5 parsecs (11.4 light years); its half-mass radius is about 3.9 parsecs (12.7 light years); and its tidal radius is about 12 parsecs (39 light years). However, the tidal radius also includes many stars that are merely "passing through" and not bona fide cluster members. Altogether, the cluster contains at least 1000 gravitationally bound stars, for a total mass of about 500–600 Solar masses. A recent survey counts 1010 high-probability members, of which 68% are M dwarfs, 30% are Sun-like stars of spectral classes F, G, and K, and about 2% are bright stars of spectral class A. Also present are five giant stars, four of which have spectral class K0 III and the fifth G0 III. So far, eleven white dwarfs have been identified, representing the final evolutionary phase of the cluster's most massive stars, which originally belonged to spectral type B. Brown dwarfs, however, are rare in this cluster, probably because they have been lost by tidal stripping from the halo. A brown dwarf has been found in the eclipsing binary system AD 3116. The cluster has a visual brightness of magnitude 3.7. Its brightest stars are blue-white and of magnitude 6 to 6.5. 42 Cancri is a confirmed member. Planets In September 2012, two planets which orbit separate stars were discovered in the Beehive Cluster. The finding was significant for being the first planets detected orbiting stars like Earth's Sun that were situated in stellar clusters. Planets had previously been detected in such clusters, but not orbiting stars like the Sun. The planets have been designated Pr0201 b and Pr0211 b. The 'b' at the end of their names indicates that the bodies are planets. The discoveries are what have been termed hot Jupiters, massive gas giants that, unlike the planet Jupiter, orbit very close to their parent stars. The announcement describing the planetary finds, written by Sam Quinn as the lead author, was published in the Astrophysical Journal Letters. Quinn's team worked with David Latham of the Harvard–Smithsonian Center for Astrophysics, utilizing the Smithsonian Astrophysical Observatory's Fred Lawrence Whipple Observatory. In 2016 additional observations found a second planet in the Pr0211 system, Pr0211 c. This made Pr0211 the first multi-planet system to be discovered in an open cluster. The Kepler space telescope, in its K2 mission, discovered planets around several more stars in the Beehive Cluster. The stars K2-95, K2-100, K2-101, K2-102, K2-103, and K2-104 host a single planet each, and K2-264 has a two-planet system.	Physical sciences	Notable star clusters	Astronomy
962171	https://en.wikipedia.org/wiki/Stern%E2%80%93Gerlach%20experiment	Stern–Gerlach experiment	In quantum physics, the Stern–Gerlach experiment demonstrated that the spatial orientation of angular momentum is quantized. Thus an atomic-scale system was shown to have intrinsically quantum properties. In the original experiment, silver atoms were sent through a spatially-varying magnetic field, which deflected them before they struck a detector screen, such as a glass slide. Particles with non-zero magnetic moment were deflected, owing to the magnetic field gradient, from a straight path. The screen revealed discrete points of accumulation, rather than a continuous distribution, owing to their quantized spin. Historically, this experiment was decisive in convincing physicists of the reality of angular-momentum quantization in all atomic-scale systems. After its conception by Otto Stern in 1921, the experiment was first successfully conducted with Walther Gerlach in early 1922. Description The Stern–Gerlach experiment involves sending silver atoms through an inhomogeneous magnetic field and observing their deflection. Silver atoms were evaporated using an electric furnace in a vacuum. Using thin slits, the atoms were guided into a flat beam and the beam sent through an inhomogeneous magnetic field before colliding with a metallic plate. The laws of classical physics predict that the collection of condensed silver atoms on the plate should form a thin solid line in the same shape as the original beam. However, the inhomogeneous magnetic field caused the beam to split in two separate directions, creating two lines on the metallic plate. The results show that particles possess an intrinsic angular momentum that is closely analogous to the angular momentum of a classically spinning object, but that takes only certain quantized values. Another important result is that only one component of a particle's spin can be measured at one time, meaning that the measurement of the spin along the z-axis destroys information about a particle's spin along the x and y axis. The experiment is normally conducted using electrically neutral particles such as silver atoms. This avoids the large deflection in the path of a charged particle moving through a magnetic field and allows spin-dependent effects to dominate. If the particle is treated as a classical spinning magnetic dipole, it will precess in a magnetic field because of the torque that the magnetic field exerts on the dipole (see torque-induced precession). If it moves through a homogeneous magnetic field, the forces exerted on opposite ends of the dipole cancel each other out and the trajectory of the particle is unaffected. However, if the magnetic field is inhomogeneous then the force on one end of the dipole will be slightly greater than the opposing force on the other end, so that there is a net force which deflects the particle's trajectory. If the particles were classical spinning objects, one would expect the distribution of their spin angular momentum vectors to be random and continuous. Each particle would be deflected by an amount proportional to the dot product of its magnetic moment with the external field gradient, producing some density distribution on the detector screen. Instead, the particles passing through the Stern–Gerlach apparatus are deflected either up or down by a specific amount. This was a measurement of the quantum observable now known as spin angular momentum, which demonstrated possible outcomes of a measurement where the observable has a discrete set of values or point spectrum. Although some discrete quantum phenomena, such as atomic spectra, were observed much earlier, the Stern–Gerlach experiment allowed scientists to directly observe separation between discrete quantum states for the first time. Theoretically, quantum angular momentum of any kind has a discrete spectrum, which is sometimes briefly expressed as "angular momentum is quantized". Experiment using particles with +1/2 or −1/2 spin If the experiment is conducted using charged particles like electrons, there will be a Lorentz force that tends to bend the trajectory in a circle. This force can be cancelled by an electric field of appropriate magnitude oriented transverse to the charged particle's path. Electrons are spin-1/2 particles. These have only two possible spin angular momentum values measured along any axis, or , a purely quantum mechanical phenomenon. Because its value is always the same, it is regarded as an intrinsic property of electrons, and is sometimes known as "intrinsic angular momentum" (to distinguish it from orbital angular momentum, which can vary and depends on the presence of other particles). If one measures the spin along a vertical axis, electrons are described as "spin up" or "spin down", based on the magnetic moment pointing up or down, respectively. To mathematically describe the experiment with spin particles, it is easiest to use Dirac's bra–ket notation. As the particles pass through the Stern–Gerlach device, they are deflected either up or down, and observed by the detector which resolves to either spin up or spin down. These are described by the angular momentum quantum number , which can take on one of the two possible allowed values, either or . The act of observing (measuring) the momentum along the axis corresponds to the -axis angular momentum operator, often denoted . In mathematical terms, the initial state of the particles is where constants and are complex numbers. This initial state spin can point in any direction. The squares of the absolute values and are respectively the probabilities for a system in the state to be found in and after the measurement along axis is made. The constants and must also be normalized in order that the probability of finding either one of the values be unity, that is we must ensure that . However, this information is not sufficient to determine the values of and , because they are complex numbers. Therefore, the measurement yields only the squared magnitudes of the constants, which are interpreted as probabilities. Sequential experiments If we link multiple Stern–Gerlach apparatuses (the rectangles containing S-G), we can clearly see that they do not act as simple selectors, i.e. filtering out particles with one of the states (pre-existing to the measurement) and blocking the others. Instead they alter the state by observing it (as in light polarization). In the figure below, x and z name the directions of the (inhomogenous) magnetic field, with the x-z-plane being orthogonal to the particle beam. In the three S-G systems shown below, the cross-hatched squares denote the blocking of a given output, i.e. each of the S-G systems with a blocker allows only particles with one of two states to enter the next S-G apparatus in the sequence. Experiment 1 The top illustration shows that when a second, identical, S-G apparatus is placed at the exit of the first apparatus, only z+ is seen in the output of the second apparatus. This result is expected since all particles at this point are expected to have z+ spin, as only the z+ beam from the first apparatus entered the second apparatus. Experiment 2 The middle system shows what happens when a different S-G apparatus is placed at the exit of the z+ beam resulting of the first apparatus, the second apparatus measuring the deflection of the beams on the x axis instead of the z axis. The second apparatus produces x+ and x- outputs. Now classically we would expect to have one beam with the x characteristic oriented + and the z characteristic oriented +, and another with the x characteristic oriented - and the z characteristic oriented +. Experiment 3 The bottom system contradicts that expectation. The output of the third apparatus which measures the deflection on the z axis again shows an output of z- as well as z+. Given that the input to the second S-G apparatus consisted only of z+, it can be inferred that a S-G apparatus must be altering the states of the particles that pass through it. This experiment can be interpreted to exhibit the uncertainty principle: since the angular momentum cannot be measured on two perpendicular directions at the same time, the measurement of the angular momentum on the x direction destroys the previous determination of the angular momentum in the z direction. That's why the third apparatus measures renewed z+ and z- beams like the x measurement really made a clean slate of the z+ output. History The Stern–Gerlach experiment was conceived by Otto Stern in 1921 and performed by him and Walther Gerlach in Frankfurt in 1922. At the time of the experiment, the most prevalent model for describing the atom was the Bohr-Sommerfeld model, which described electrons as going around the positively charged nucleus only in certain discrete atomic orbitals or energy levels. Since the electron was quantized to be only in certain positions in space, the separation into distinct orbits was referred to as space quantization. The Stern–Gerlach experiment was meant to test the Bohr–Sommerfeld hypothesis that the direction of the angular momentum of a silver atom is quantized. The experiment was first performed with an electromagnet that allowed the non-uniform magnetic field to be turned on gradually from a null value. When the field was null, the silver atoms were deposited as a single band on the detecting glass slide. When the field was made stronger, the middle of the band began to widen and eventually to split into two, so that the glass-slide image looked like a lip-print, with an opening in the middle, and closure at either end. In the middle, where the magnetic field was strong enough to split the beam into two, statistically half of the silver atoms had been deflected by the non-uniformity of the field. Note that the experiment was performed several years before George Uhlenbeck and Samuel Goudsmit formulated their hypothesis about the existence of electron spin in 1925. Even though the result of the Stern−Gerlach experiment has later turned out to be in agreement with the predictions of quantum mechanics for a spin-1/2 particle, the experimental result was also consistent with the Bohr–Sommerfeld theory. In 1927, T.E. Phipps and J.B. Taylor reproduced the effect using hydrogen atoms in their ground state, thereby eliminating any doubts that may have been caused by the use of silver atoms. However, in 1926 the non-relativistic scalar Schrödinger equation had incorrectly predicted the magnetic moment of hydrogen to be zero in its ground state. To correct this problem Wolfgang Pauli considered a spin-1/2 version of the Schrödinger equation using the 3 Pauli matrices which now bear his name, which was later shown by Paul Dirac in 1928 to be a consequence of his relativistic Dirac equation. In the early 1930s Stern, together with Otto Robert Frisch and Immanuel Estermann improved the molecular beam apparatus sufficiently to measure the magnetic moment of the proton, a value nearly 2000 times smaller than the electron moment. In 1931, theoretical analysis by Gregory Breit and Isidor Isaac Rabi showed that this apparatus could be used to measure nuclear spin whenever the electronic configuration of the atom was known. The concept was applied by Rabi and Victor W. Cohen in 1934 to determine the spin of sodium atoms. In 1938 Rabi and coworkers inserted an oscillating magnetic field element into their apparatus, inventing nuclear magnetic resonance spectroscopy. By tuning the frequency of the oscillator to the frequency of the nuclear precessions they could selectively tune into each quantum level of the material under study. Rabi was awarded the Nobel Prize in 1944 for this work. Importance The Stern–Gerlach experiment was the first direct evidence of angular-momentum quantization in quantum mechanics, and it strongly influenced later developments in modern physics: In the decade that followed, scientists showed using similar techniques, that the nuclei of some atoms also have quantized angular momentum. It is the interaction of this nuclear angular momentum with the spin of the electron that is responsible for the hyperfine structure of the spectroscopic lines. Norman F. Ramsey later modified the Rabi apparatus to improve its sensitivity (using the separated oscillatory field method). In the early sixties, Ramsey, H. Mark Goldenberg, and Daniel Kleppner used a Stern–Gerlach system to produce a beam of polarized hydrogen as the source of energy for the hydrogen maser. This led to developing an extremely stable clock based on a hydrogen maser. From 1967 until 2019, the second was defined based on 9,192,631,770 Hz hyperfine transition of a cesium-133 atom; the atomic clock which is used to set this standard is an application of Ramsey's work. The Stern–Gerlach experiment has become a prototype for quantum measurement, demonstrating the observation of a discrete value (eigenvalue) of a physical property, previously assumed to be continuous. Entering the Stern–Gerlach magnet, the direction of the silver atom's magnetic moment is indefinite, but when the atom is registered at the screen, it is observed to be at either one spot or the other, and this outcome cannot be predicted in advance. Because the experiment illustrates the character of quantum measurements, The Feynman Lectures on Physics use idealized Stern–Gerlach apparatuses to explain the basic mathematics of quantum theory.	Physical sciences	Quantum mechanics	Physics
962527	https://en.wikipedia.org/wiki/Viviparous%20lizard	Viviparous lizard	The viviparous lizard, or common lizard, (Zootoca vivipara, formerly Lacerta vivipara) is a Eurasian lizard. It lives farther north than any other species of non-marine reptile, and is named for the fact that it is viviparous, meaning it gives birth to live young (although they will sometimes lay eggs normally). Both "Zootoca" and "vivipara" mean "live birth", in (Latinized) Greek and Latin respectively. It was called Lacerta vivipara until the genus Lacerta was split into nine genera in 2007 by Arnold, Arribas & Carranza. Male and female Zootoca vivipara are equally likely to contract blood parasites. Additionally, larger males have been shown to reproduce more times in a given reproductive season than smaller ones. The lizard is also unique as it is exclusively carnivorous, eating only flies, spiders, and insects. Studies show that the more carnivorous an individual is (the more insects they eat), the less diverse the population of parasitic helminths that infest the lizards. Zootoca vivipara lives in very cold climates, yet participates in normal thermoregulation instead of thermoconformity. They have the largest range of all terrestrial lizards which even include subarctic regions. It is able to survive these harsh climates as individuals will freeze in especially cold seasons and thaw two months later. They also live closer to geological phenomena that provide a warmer environment for them. Description Zootoca vivipara is a small lizard, with an average length between {150-200 mm} . They exhibit no particular colour, but can be brown, red, grey, green, or black. The species exhibits some sexual dimorphisms. Female Z. vivipara undergo colour polymorphism more commonly than males. A female lizard's display differs in ventral coloration, ranging from pale yellow to bright orange and a mixed coloration. There have been many hypotheses for the genetic cause of this polymorphic coloration. These hypothesis test for coloration due to thermoregulation, predator avoidance, and social cues, specifically sexual reproduction. Through an experiment conducted by Vercken et al., colour polymorphism in viviparous lizard is caused by social cues, rather than the other hypotheses. More specifically, the ventral coloration that is seen in female lizards is associated with patterns of sexual reproduction and sex allocation. The underside of the male is typically more colourful and bright, with yellow, orange, green, and blue, and the male typically has spots along its back. On the other hand, females typically have darker stripes down their backs and sides. Additionally, males have been found to have larger heads than their female counterparts, and this trait appears to be sexually selected for. Males with larger heads are more likely to be successful in mating and male-male interactions than smaller-headed Z. vivipara. Larger males also have been shown to reproduce more frequently during one mating season compared to smaller males. Characteristic behaviors of the species includes tongue flicking in the presence of a predator and female-female aggression that seems to be mediated by the colour of their side stripe. Habitat and distribution Habitat Z. vivipara is terrestrial, so they spend most of their time on the ground, though they do occasionally visit sites of higher elevation. The lizard thermoregulates by basking in the sun for much of the time. In colder weather, they have been known to hibernate to maintain proper body temperatures. They hibernate between October and March. Their typical habitats include heathland, moorland, woodland and grassland. The viviparous lizard is native to much of northern Eurasia. In Europe, it is mainly found north of the Alps and the Carpathians, including the British Isles but not Iceland, as well as in parts of northern Iberia and the Balkans; In Asia it is mostly found in Russia, excluding northern Siberia, and in northern Kazakhstan, Mongolia, China, and Japan. Z. vivipara has the largest distribution of any species of lizard in the world. Home range The size of the home range of the lizard ranges from 539 m2 to 1692 m2, with males generally having larger home ranges. The size of an individual lizard's home range is also dependent on population density and the presence of prey. Ecology Diet Unlike many other lizards, Z. vivipara is exclusively carnivorous. Their diet consists of flies, spiders, and various other insects, including hemipterans (such as cicadas), moth larvae, and mealworms. The species is a predator, so it actively hunts down all of its prey. One study found that when controlled for body size, females consumed more food than males. Feeding rates also increased with increased sunshine. Predation Birds are common predators of Z. vivipara. Male-biased predation of Z. vivipara by the great grey shrike (L. excubitor) has been studied, finding that adult males, over adult females and juveniles, were preferentially predated on. This bias may be due to increased activity of adult males during the reproductive season. Predators of this species include birds of prey, crows, snakes, shrikes, hedgehogs, shrews, foxes, and domestic cats. Diseases and parasites Z. vivpara can be infested by helminths, a small parasitic worm. The species diversity of parasites is affected by the diet of the individual lizard and the number of parasites on a host is affected by the host's size. Results of a study shows that the more carnivorous an individual is, the less diverse its parasite population. Additionally, larger lizards had a greater number of parasites on them. Z. vivipara is also infected by blood parasites. In a study investigating the prevalence of blood parasites in Z. vivipara and L. agilis, Z. vivipara was found to be parasitized with an incidence rate of 39.8%, while L. agilis was parasitized with an incidence rate of 22.3%. This same study shows that there was not a significant difference between the parasitization of male and female Z. vivipara. Reproduction and life history Viviparity and oviparity The viviparous lizard is named as such because it is viviparous. This refers to its ability to give birth to live young, although the lizards are also able to lay eggs. The origin of this characteristic is under debate. Some scientists argue that viviparity evolved from oviparity, or the laying of eggs, only once. Proponents of this theory also argue that if this is the case, it is possible, though rare, for species to transition back to oviparity. Research from Yann Surget-Groba suggests that there have in fact been multiple events of the evolution of viviparity from oviparity across different clades of the viviparous lizard. They also argue that a reversion to oviparity is not as rare as once believed, but has occurred 2 to 3 times in the history of the species. The range of viviparous populations of Z. vivipara extends from France to Russia. Oviparous populations are only found in northern Spain and the southwest of France. Some research in the Italian alps has suggested that distinct populations of oviparous and viviparous Z. vivipara should be considered separate species. Cornetti et al. (2015) identified that viviparous and oviparous subpopulations in contact with each other in the Italian alps are reproductively isolated. Hybridization between viviparous and oviparous individuals of Z. vivipara leads to embryonic malformations in the laboratory. However, these crosses do produce a "hybridized" generation of offspring, with females retaining embryos for much longer in utero than oviparous females, with embryos surrounded by thin, translucent shells. Fertilization Z. vivipara juveniles reach sexual maturity during their second year of their life. A study that explored the presence of male sex cells in reproducing males found that for the two weeks following the end of hibernation, males are infertile, and therefore incapable of reproducing. The same study also found that larger males produce more sperm during the reproductive season and have fewer left over at the end of the reproductive season than their smaller counterparts. This suggests that the larger a male is, the more reproductive events they participate in. Brood size Research also suggests that in exclusively oviparous populations of Z. vivipara, altitude influences the number of clutches laid in a reproductive season as well as when reproduction begins. Generally, lizards living at higher altitudes have been found to begin reproduction later and lay fewer clutches (often 1) in a given reproductive season. Life span Z. vivipara typically lives for 5 to 6 years. Mating Mate searching behavior Head size is a sexually dimorphic trait, with males having larger heads than females. The average head width and length of the males measured were found to be , respectively. The average head width and length of the females measured were found to be , respectively. During the first state of courtship in Z. vivipara, called "Capture", the male uses its mouth and jaw to capture the female and initiate copulation. The results of this study demonstrated that males with larger head sizes (both length and width) were more successful in mating than those with smaller heads, suggesting that head size undergoes sexual selection. Male-male interaction Head size has also been shown to be a predictor of success in male-male interactions. The head is used as a weapon in male-male interactions, and a larger head is typically more effective, leading to greater success during male-male aggressive encounters. This aggression and interaction is centered around available mates, so males with smaller heads have significantly less access to females for reproduction. Thermoregulation This lizard has an exceptionally large range that includes subarctic geography. As a result, thermoregulation is necessary for the thermal homeostasis of the species. Typically, in temperature extremes, a species will adopt the behavioral strategy of thermoconformity, where they do not actively thermoregulate, but adapt to survive in the harsh temperature. This occurs because the cost of thermoregulating in such an extreme environment becomes too high and begins to outweigh the benefits. Despite this, Z. vivipara still employs the strategy of thermoregulation, like basking. Thermoregulation is important in Z. vivipara as it allows for proper locomotive performance, escape behavior, and other key behaviors for survival. The ability of Z. vivipara to thermoregulate in such harsh environments has been attributed to two primary reasons. The first is that Z. vivipara has remarkable behaviors to combat the cold, and there are geological phenomena in their distribution that maintains their habitats at a temperature that the species can survive in. One of the specific behaviors used to combat the extreme cold is a "supercooled" state. Z. vivipara remains in this state through the winter until temperatures dropped below . After that, individuals completely froze until they were thawed by warmer weather later in the year, often 2 months later. Despite very cold air in the subarctic habitats of these lizards, the soil-heating effects of unfrozen groundwater has been observed regulating the temperature of their soil habitats. They find warm microhabitats that do not drop below the freezing point of their body fluids. These lizards have exceptional hardiness to the cold, which allows them to hibernate in upper soil layers in temperatures as low as . This cold hardiness along with the favorable hydrogeological conditions of groundwater-warmed soil habitats allows for the wide distribution of lizards throughout the palearctic. Colour polymorphism The colour polymorphism of female Z. vivipara has not been thoroughly studied in past years, regardless of the extensive research done on the species itself. Females exhibit three types of body colouration within a population: yellow, orange, and mixture of the two. These discrete traits are inherited maternally and exist throughout the individual's lifetime. The organism's colour morphs are determined by their genotype as well as their environment. The frequency of multiple morphs occurring in a population varies with the level of population density and frequency-dependent environments. These factors cause the lizards to vary in terms of their fitness (clutch size, sex ratio, hatching success). In lower density populations, colour polymorphism is more prevalent. This is because viviparous lizards thrive in environments where intraspecific competition is low. Increased competition among individuals results in lower survival rates of lizards. Additionally, female lizards disperse through habitats based on the frequency of colour types that are already present in the population. Their reproductive abilities vary according to this frequency-dependent environment. The number of offspring that they produce correlates with the colour morph: yellow females produce the fewest offspring, while orange females produce more than yellow, but fewer than mixed females, which produce the most offspring. The amount of offspring produced varies in regards to colour frequencies in the population; for example, if yellow females have higher density within the population, the clutch size for orange lizards is usually lower. Orange females are more sensitive to intraspecific and colour-specific competition. They have smaller clutch sizes when the density of the population is high, or when the number of yellow females in the population is high. This could be due to their need to conserve energy for survival and reproductive events. Their colour morph remains in the population due to the trade-off between the size of offspring and the clutch size. Offspring born in smaller clutches are often larger and thus have a higher survival likelihood. Natural selection will favor individuals with larger size because of their advantage in physical competition with others. Yellow females have larger clutch sizes early in their life, but their hatch success decreases as the female ages. Their reproductive viability decreases, resulting in fewer offspring throughout their lifetime. Yellow morphs remain in the population due to their large clutch size, which causes an increased frequency of those females. Selection favors the yellow morph because of the ability to produce large clutch sizes, which increases the female's fitness. In mixed-coloured females, reproductive success is less sensitive to competition and frequency-dependent environments. Since these lizards show a mixture of yellow and orange colouration, they adopt benefits from both of the morphs. As a result, they can maintain high reproductive success and hatching success with large clutch sizes. Their colour morph remains in the population due to its high fitness, which selection will favor. All three colours have evolutionary advantages in different ways. While yellow females have higher fitness due to their large clutch sizes, orange females enjoy high fitness due to their large body size and increased competitive advantages. Mixed females exhibit both of these advantages.	Biology and health sciences	Lizards and other Squamata	Animals
558397	https://en.wikipedia.org/wiki/Rogue%20planet	Rogue planet	A rogue planet, also termed a free-floating planet (FFP) or an isolated planetary-mass object (iPMO), is an interstellar object of planetary mass which is not gravitationally bound to any star or brown dwarf. Rogue planets may originate from planetary systems in which they are formed and later ejected, or they can also form on their own, outside a planetary system. The Milky Way alone may have billions to trillions of rogue planets, a range the upcoming Nancy Grace Roman Space Telescope is expected to refine. Some planetary-mass objects may have formed in a similar way to stars, and the International Astronomical Union has proposed that such objects be called sub-brown dwarfs. A possible example is Cha 110913−773444, which may either have been ejected and become a rogue planet or formed on its own to become a sub-brown dwarf. Terminology The two first discovery papers use the names isolated planetary-mass objects (iPMO) and free-floating planets (FFP). Most astronomical papers use one of these terms. The term rogue planet is more often used for microlensing studies, which also often uses the term FFP. A press release intended for the public might use an alternative name. The discovery of at least 70 FFPs in 2021, for example, used the terms rogue planet, starless planet, wandering planet and free-floating planet in different press releases. Discovery Isolated planetary-mass objects (iPMO) were first discovered in 2000 by the UK team Lucas & Roche with UKIRT in the Orion Nebula. In the same year the Spanish team Zapatero Osorio et al. discovered iPMOs with Keck spectroscopy in the σ Orionis cluster. The spectroscopy of the objects in the Orion Nebula was published in 2001. Both European teams are now recognized for their quasi-simultaneous discoveries. In 1999 the Japanese team Oasa et al. discovered objects in Chamaeleon I that were spectroscopically confirmed years later in 2004 by the US team Luhman et al. Observation There are two techniques to discover free-floating planets: direct imaging and microlensing. Microlensing Astrophysicist Takahiro Sumi of Osaka University in Japan and colleagues, who form the Microlensing Observations in Astrophysics and the Optical Gravitational Lensing Experiment collaborations, published their study of microlensing in 2011. They observed 50 million stars in the Milky Way by using the MOA-II telescope at New Zealand's Mount John Observatory and the University of Warsaw telescope at Chile's Las Campanas Observatory. They found 474 incidents of microlensing, ten of which were brief enough to be planets of around Jupiter's size with no associated star in the immediate vicinity. The researchers estimated from their observations that there are nearly two Jupiter-mass rogue planets for every star in the Milky Way. One study suggested a much larger number, up to 100,000 times more rogue planets than stars in the Milky Way, though this study encompassed hypothetical objects much smaller than Jupiter. A 2017 study by Przemek Mróz of Warsaw University Observatory and colleagues, with six times larger statistics than the 2011 study, indicates an upper limit on Jupiter-mass free-floating or wide-orbit planets of 0.25 planets per main-sequence star in the Milky Way. In September 2020, astronomers using microlensing techniques reported the detection, for the first time, of an Earth-mass rogue planet (named OGLE-2016-BLG-1928) unbound to any star and free floating in the Milky Way galaxy. Direct imaging Microlensing planets can only be studied by the microlensing event, which makes the characterization of the planet difficult. Astronomers therefore turn to isolated planetary-mass objects (iPMO) that were found via the direct imaging method. To determine a mass of a brown dwarf or iPMO one needs for example the luminosity and the age of an object. Determining the age of a low-mass object has proven to be difficult. It is no surprise that the vast majority of iPMOs are found inside young nearby star-forming regions of which astronomers know their age. These objects are younger than 200 Myrs, are massive (>5 ) and belong to the L- and T-dwarfs. There is however a small growing sample of cold and old Y-dwarfs that have estimated masses of 8-20 . Nearby rogue planet candidates of spectral type Y include WISE 0855−0714 at a distance of . If this sample of Y-dwarfs can be characterized with more accurate measurements or if a way to better characterize their ages can be found, the number of old and cold iPMOs will likely increase significantly. The first iPMOs were discovered in the early 2000s via direct imaging inside young star-forming regions. These iPMOs found via direct imaging formed probably like stars (sometimes called sub-brown dwarf). There might be iPMOs that form like a planet, which are then ejected. These objects will however be kinematically different from their natal star-forming region, should not be surrounded by a circumstellar disk and have high metallicity. None of the iPMOs found inside young star-forming regions show a high velocity compared to their star-forming region. For old iPMOs the cold WISE J0830+2837 shows a Vtan of about 100 km/s, which is high, but still consistent with formation in our galaxy. For WISE 1534–1043 one alternative scenario explains this object as an ejected exoplanet due to its high Vtan of about 200 km/s, but its color suggests it is an old metal-poor brown dwarf. Most astronomers studying massive iPMOs believe that they represent the low-mass end of the star-formation process. Astronomers have used the Herschel Space Observatory and the Very Large Telescope to observe a very young free-floating planetary-mass object, OTS 44, and demonstrate that the processes characterizing the canonical star-like mode of formation apply to isolated objects down to a few Jupiter masses. Herschel far-infrared observations have shown that OTS 44 is surrounded by a disk of at least 10 Earth masses and thus could eventually form a mini planetary system. Spectroscopic observations of OTS 44 with the SINFONI spectrograph at the Very Large Telescope have revealed that the disk is actively accreting matter, similar to the disks of young stars. Binaries The first discovery of a resolved planetary-mass binary was 2MASS J1119–1137AB. There are however other binaries known, such as 2MASS J1553022+153236AB, WISE 1828+2650, WISE 0146+4234, WISE J0336−0143 (could also be a brown dwarf and a planetary-mass object (BD+PMO) binary), NIRISS-NGC1333-12 and several objects discovered by Zhang et al. In the Orion Nebula a population of 40 wide binaries and 2 triple systems were discovered. This was surprising for two reasons: The trend of binaries of brown dwarfs predicted a decrease of distance between low mass objects with decreasing mass. It was also predicted that the binary fraction decreases with mass. These binaries were named Jupiter-mass binary objects (JuMBOs). They make up at least 9% of the iPMOs and have a separation smaller than 340 AU. It is unclear how these JuMBOs formed, but an extensive study argued that they formed in situ, like stars. If they formed like stars, then there must be an unknown "extra ingredient" to allow them to form. If they formed like planets and were later ejected, then it has to be explained why these binaries did not break apart during the ejection process. Future measurements with JWST might resolve if these objects formed as ejected planets or as stars. A study by Kevin Luhman reanalysed the NIRCam data and found that most JuMBOs did not appear in his sample of substellar objects. Moreover the color were consistent with reddened background sources or low signal-to-noise sources. Only JuMBO 29 is identified as a good candidate in this work. JuMBO 29 also was observed with NIRSpec and one component was identified as a young M8 source. This spectral type is consistent with a low mass for the age of the Orion Nebula. Total number of known iPMOs There are likely hundreds of known candidate iPMOs, over a hundred objects with spectra and a small but growing number of candidates discovered via microlensing. Some large surveys include: As of December 2021, the largest-ever group of rogue planets was discovered, numbering at least 70 and up to 170 depending on the assumed age. They are found in the OB association between Upper Scorpius and Ophiuchus with masses between 4 and 13 and age around 3 to 10 million years, and were most likely formed by either gravitational collapse of gas clouds, or formation in a protoplanetary disk followed by ejection due to dynamical instabilities. Follow-up observations with spectroscopy from the Subaru Telescope and Gran Telescopio Canarias showed that the contamination of this sample is quite low (≤6%). The 16 young objects had a mass between 3 and 14 , confirming that they are indeed planetary-mass objects. In October 2023 an even larger group of 540 planetary-mass object candidates was discovered in the Trapezium Cluster and inner Orion Nebula with JWST. The objects have a mass between 13 and 0.6 . A surprising number of these objects formed wide binaries, which was not predicted. Formation There are in general two scenarios that can lead to the formation of an isolated planetary-mass object (iPMO). It can form like a planet around a star and is then ejected, or it forms like a low-mass star or brown dwarf in isolation. This can influence its composition and motion. Formation like a star Objects with a mass of at least one Jupiter mass were thought to be able to form via collapse and fragmentation of molecular clouds from models in 2001. Pre-JWST observations have shown that objects below 3-5 are unlikely to form on their own. Observations in 2023 in the Trapezium Cluster with JWST have shown that objects as massive as 0.6 might form on their own, not requiring a steep cut-off mass. A particular type of globule, called globulettes, are thought to be birthplaces for brown dwarfs and planetary-mass objects. Globulettes are found in the Rosette Nebula and IC 1805. Sometimes young iPMOs are still surrounded by a disk that could form exomoons. Due to the tight orbit of this type of exomoon around their host planet, they have a high chance of 10-15% to be transiting. Disks Some very young star-forming regions, typically younger than 5 million years, sometimes contain isolated planetary-mass objects with infrared excess and signs of accretion. Most well known is the iPMO OTS 44 discovered to have a disk and being located in Chamaeleon I. Charmaeleon I and II have other candidate iPMOs with disks. Other star-forming regions with iPMOs with disks or accretion are Lupus I, Rho Ophiuchi Cloud Complex, Sigma Orionis cluster, Orion Nebula, Taurus, NGC 1333 and IC 348. A large survey of disks around brown dwarfs and iPMOs with ALMA found that these disks are not massive enough to form earth-mass planets. There is still the possibility that the disks already have formed planets. Studies of red dwarfs have shown that some have gas-rich disks at a relative old age. These disks were dubbed Peter Pan Disks and this trend could continue into the planetary-mass regime. One Peter Pan disk is the 45 Myr old brown dwarf 2MASS J02265658-5327032 with a mass of about 13.7 , which is close to the planetary-mass regime. Recent studies of the nearby planetary-mass object 2MASS J11151597+1937266 found that this nearby iPMO is surrounded by a disk. It shows signs of accretion from the disk and also infrared excess. Formation like a planet Ejected planets are predicted to be mostly low-mass (<30 Figure 1 Ma et al.) and their mean mass depends on the mass of their host star. Simulations by Ma et al. did show that 17.5% of 1 stars eject a total of 16.8 per star with a typical (median) mass of 0.8 for an individual free-floating planet (FFP). For lower mass red dwarfs with a mass of 0.3 12% of stars eject a total of 5.1 per star with a typical mass of 0.3 for an individual FFP. Hong et al. predicted that exomoons can be scattered by planet-planet interactions and become ejected exomoons. Higher mass (0.3-1 ) ejected FFP are predicted to be possible, but they are also predicted to be rare. Ejection of a planet can occur via planet-planet scatter or due a stellar flyby. Another possibility is the ejection of a fragment of a disk that then forms into a planetary-mass object. Another suggested scenario is the ejection of planets in a tilted circumbinary orbit. Interactions with the central binary and the planets with each other can lead to the ejection of the lower-mass planet in the system. Other scenarios If a stellar or brown dwarf embryo experiences a halted accretion, it could remain low-mass enough to become a planetary-mass object. Such a halted accretion could occur if the embryo is ejected or if its circumstellar disk experiences photoevaporation near O-stars. Objects that formed via the ejected embryo scenario would have smaller or no disk and the fraction of binaries decreases for such objects. It could also be that free-floating planetary-mass objects for from a combination of scenarios. Fate Most isolated planetary-mass objects will float in interstellar space forever. Some iPMOs will have a close encounter with a planetary system. This rare encounter can have three outcomes: The iPMO will remain unbound, it could be weakly bound to the star, or it could "kick out" the exoplanet, replacing it. Simulations have shown that the vast majority of these encounters result in a capture event with the iPMO being weakly bound with a low gravitational binding energy and an elongated highly eccentric orbit. These orbits are not stable and 90% of these objects gain energy due to planet-planet encounters and are ejected back into interstellar space. Only 1% of all stars will experience this temporary capture. Warmth Interstellar planets generate little heat and are not heated by a star. However, in 1998, David J. Stevenson theorized that some planet-sized objects adrift in interstellar space might sustain a thick atmosphere that would not freeze out. He proposed that these atmospheres would be preserved by the pressure-induced far-infrared radiation opacity of a thick hydrogen-containing atmosphere. During planetary-system formation, several small protoplanetary bodies may be ejected from the system. An ejected body would receive less of the stellar-generated ultraviolet light that can strip away the lighter elements of its atmosphere. Even an Earth-sized body would have enough gravity to prevent the escape of the hydrogen and helium in its atmosphere. In an Earth-sized object the geothermal energy from residual core radioisotope decay could maintain a surface temperature above the melting point of water, allowing liquid-water oceans to exist. These planets are likely to remain geologically active for long periods. If they have geodynamo-created protective magnetospheres and sea floor volcanism, hydrothermal vents could provide energy for life. These bodies would be difficult to detect because of their weak thermal microwave radiation emissions, although reflected solar radiation and far-infrared thermal emissions may be detectable from an object that is less than 1,000 astronomical units from Earth. Around five percent of Earth-sized ejected planets with Moon-sized natural satellites would retain their satellites after ejection. A large satellite would be a source of significant geological tidal heating. List The table below lists rogue planets, confirmed or suspected, that have been discovered. It is yet unknown whether these planets were ejected from orbiting a star or else formed on their own as sub-brown dwarfs. Whether exceptionally low-mass rogue planets (such as OGLE-2012-BLG-1323 and KMT-2019-BLG-2073) are even capable of being formed on their own is currently unknown. Discovered via direct imaging These objects were discovered with the direct imaging method. Many were discovered in young star-clusters or stellar associations and a few old are known (such as WISE 0855−0714). List is sorted after discovery year. Discovered via microlensing These objects were discovered via microlensing. Rogue planets discovered via microlensing can only be studied by the lensing event. Some of them could also be exoplanets in a wide orbit around an unseen star. Discovered via transit	Physical sciences	Planetary science	Astronomy
558685	https://en.wikipedia.org/wiki/Natural%20environment	Natural environment	The natural environment or natural world encompasses all biotic and abiotic things occurring naturally, meaning in this case not artificial. The term is most often applied to Earth or some parts of Earth. This environment encompasses the interaction of all living species, climate, weather and natural resources that affect human survival and economic activity. The concept of the natural environment can be distinguished as components: Complete ecological units that function as natural systems without massive civilized human intervention, including all vegetation, microorganisms, soil, rocks, plateaus, mountains, the atmosphere and natural phenomena that occur within their boundaries and their nature. Universal natural resources and physical phenomena that lack clear-cut boundaries, such as air, water and climate, as well as energy, radiation, electric charge and magnetism, not originating from civilized human actions. In contrast to the natural environment is the built environment. Built environments are where humans have fundamentally transformed landscapes such as urban settings and agricultural land conversion, the natural environment is greatly changed into a simplified human environment. Even acts which seem less extreme, such as building a mud hut or a photovoltaic system in the desert, the modified environment becomes an artificial one. Though many animals build things to provide a better environment for themselves, they are not human, hence beaver dams and the works of mound-building termites are thought of as natural. People cannot find absolutely natural environments on Earth,naturalness usually varies in a continuum, from 100% natural in one extreme to 0% natural in the other. The massive environmental changes of humanity in the Anthropocene have fundamentally effected all natural environments including: climate change, biodiversity loss and pollution from plastic and other chemicals in the air and water. More precisely, we can consider the different aspects or components of an environment, and see that their degree of naturalness is not uniform. If, for instance, in an agricultural field, the mineralogic composition and the structure of its soil are similar to those of an undisturbed forest soil, but the structure is quite different. Composition Earth science generally recognizes four spheres, the lithosphere, the hydrosphere, the atmosphere and the biosphere as correspondent to rocks, water, air and life respectively. Some scientists include as part of the spheres of the Earth, the cryosphere (corresponding to ice) as a distinct portion of the hydrosphere, as well as the pedosphere (to soil) as an active and intermixed sphere. Earth science (also known as geoscience, the geographical sciences or the Earth Sciences), is an all-embracing term for the sciences related to the planet Earth. There are four major disciplines in earth sciences, namely geography, geology, geophysics and geodesy. These major disciplines use physics, chemistry, biology, chronology and mathematics to build a qualitative and quantitative understanding of the principal areas or spheres of Earth. Geological activity The Earth's crust or lithosphere, is the outermost solid surface of the planet and is chemically, physically and mechanically different from underlying mantle. It has been generated greatly by igneous processes in which magma cools and solidifies to form solid rock. Beneath the lithosphere lies the mantle which is heated by the decay of radioactive elements. The mantle though solid is in a state of rheic convection. This convection process causes the lithospheric plates to move, albeit slowly. The resulting process is known as plate tectonics. Volcanoes result primarily from the melting of subducted crust material or of rising mantle at mid-ocean ridges and mantle plumes. Water on Earth Most water is found in various kinds of natural body of water. Oceans An ocean is a major body of saline water and a component of the hydrosphere. Approximately 71% of the surface of the Earth (an area of some 362 million square kilometers) is covered by ocean, a continuous body of water that is customarily divided into several principal oceans and smaller seas. More than half of this area is over 3,000 meters (9,800 ft) deep. Average oceanic salinity is around 35 parts per thousand (ppt) (3.5%), and nearly all seawater has a salinity in the range of 30 to 38 ppt. Though generally recognized as several separate oceans, these waters comprise one global, interconnected body of salt water often referred to as the World Ocean or global ocean. The deep seabeds are more than half the Earth's surface, and are among the least-modified natural environments. The major oceanic divisions are defined in part by the continents, various archipelagos and other criteria, these divisions are : (in descending order of size) the Pacific Ocean, the Atlantic Ocean, the Indian Ocean, the Southern Ocean and the Arctic Ocean. Rivers A river is a natural watercourse, usually freshwater, flowing toward an ocean, a lake, a sea or another river. A few rivers simply flow into the ground and dry up completely without reaching another body of water. The water in a river is usually in a channel, made up of a stream bed between banks. In larger rivers there is often also a wider floodplain shaped by waters over-topping the channel. Flood plains may be very wide in relation to the size of the river channel. Rivers are a part of the hydrological cycle. Water within a river is generally collected from precipitation through surface runoff, groundwater recharge, springs and the release of water stored in glaciers and snowpacks. Small rivers may also be called by several other names, including stream, creek and brook. Their current is confined within a bed and stream banks. Streams play an important corridor role in connecting fragmented habitats and thus in conserving biodiversity. The study of streams and waterways in general is known as surface hydrology. Lakes A lake (from Latin lacus) is a terrain feature, a body of water that is localized to the bottom of basin. A body of water is considered a lake when it is inland, is not part of an ocean and is larger and deeper than a pond. Natural lakes on Earth are generally found in mountainous areas, rift zones and areas with ongoing or recent glaciation. Other lakes are found in endorheic basins or along the courses of mature rivers. In some parts of the world, there are many lakes because of chaotic drainage patterns left over from the last ice age. All lakes are temporary over geologic time scales, as they will slowly fill in with sediments or spill out of the basin containing them. Ponds A pond is a body of standing water, either natural or human-made, that is usually smaller than a lake. A wide variety of human-made bodies of water are classified as ponds, including water gardens designed for aesthetic ornamentation, fish ponds designed for commercial fish breeding and solar ponds designed to store thermal energy. Ponds and lakes are distinguished from streams by their current speed. While currents in streams are easily observed, ponds and lakes possess thermally driven micro-currents and moderate wind-driven currents. These features distinguish a pond from many other aquatic terrain features, such as stream pools and tide pools. Human impact on water Humans impact the water in different ways such as modifying rivers (through dams and stream channelization), urbanization and deforestation. These impact lake levels, groundwater conditions, water pollution, thermal pollution, and marine pollution. Humans modify rivers by using direct channel manipulation. We build dams and reservoirs and manipulate the direction of the rivers and water path. Dams can usefully create reservoirs and hydroelectric power. However, reservoirs and dams may negatively impact the environment and wildlife. Dams stop fish migration and the movement of organisms downstream. Urbanization affects the environment because of deforestation and changing lake levels, groundwater conditions, etc. Deforestation and urbanization go hand in hand. Deforestation may cause flooding, declining stream flow and changes in riverside vegetation. The changing vegetation occurs because when trees cannot get adequate water they start to deteriorate, leading to a decreased food supply for the wildlife in an area. Atmosphere, climate and weather The atmosphere of the Earth serves as a key factor in sustaining the planetary ecosystem. The thin layer of gases that envelops the Earth is held in place by the planet's gravity. Dry air consists of 78% nitrogen, 21% oxygen, 1% argon, inert gases and carbon dioxide. The remaining gases are often referred to as trace gases. The atmosphere includes greenhouse gases such as carbon dioxide, methane, nitrous oxide and ozone. Filtered air includes trace amounts of many other chemical compounds. Air also contains a variable amount of water vapor and suspensions of water droplets and ice crystals seen as clouds. Many natural substances may be present in tiny amounts in an unfiltered air sample, including dust, pollen and spores, sea spray, volcanic ash and meteoroids. Various industrial pollutants also may be present, such as chlorine (elementary or in compounds), fluorine compounds, elemental mercury, and sulphur compounds such as sulphur dioxide (SO2). The ozone layer of the Earth's atmosphere plays an important role in reducing the amount of ultraviolet (UV) radiation that reaches the surface. As DNA is readily damaged by UV light, this serves to protect life at the surface. The atmosphere also retains heat during the night, thereby reducing the daily temperature extremes. Layers of the atmosphere Principal layers Earth's atmosphere can be divided into five main layers. These layers are mainly determined by whether temperature increases or decreases with altitude. From highest to lowest, these layers are: Exosphere: The outermost layer of Earth's atmosphere extends from the exobase upward, mainly composed of hydrogen and helium. Thermosphere: The top of the thermosphere is the bottom of the exosphere, called the exobase. Its height varies with solar activity and ranges from about . The International Space Station orbits in this layer, between . In another way, the thermosphere is Earth's second highest atmospheric layer, extending from approximately 260,000 feet at the mesopause to the thermopause at altitudes ranging from 1,600,000 to 3,300,000 feet. Mesosphere: The mesosphere extends from the stratopause to . It is the layer where most meteors burn up upon entering the atmosphere. Stratosphere: The stratosphere extends from the tropopause to about . The stratopause, which is the boundary between the stratosphere and mesosphere, typically is at . Troposphere: The troposphere begins at the surface and extends to between at the poles and at the equator, with some variation due to weather. The troposphere is mostly heated by transfer of energy from the surface, so on average the lowest part of the troposphere is warmest and temperature decreases with altitude. The tropopause is the boundary between the troposphere and stratosphere. Other layers Within the five principal layers determined by temperature there are several layers determined by other properties. The ozone layer is contained within the stratosphere. It is mainly located in the lower portion of the stratosphere from about , though the thickness varies seasonally and geographically. About 90% of the ozone in our atmosphere is contained in the stratosphere. The ionosphere: The part of the atmosphere that is ionized by solar radiation, stretches from and typically overlaps both the exosphere and the thermosphere. It forms the inner edge of the magnetosphere. The homosphere and heterosphere: The homosphere includes the troposphere, stratosphere and mesosphere. The upper part of the heterosphere is composed almost completely of hydrogen, the lightest element. The planetary boundary layer is the part of the troposphere that is nearest the Earth's surface and is directly affected by it, mainly through turbulent diffusion. Effects of global warming The dangers of global warming are being increasingly studied by a wide global consortium of scientists. These scientists are increasingly concerned about the potential long-term effects of global warming on our natural environment and on the planet. Of particular concern is how climate change and global warming caused by anthropogenic, or human-made releases of greenhouse gases, most notably carbon dioxide, can act interactively and have adverse effects upon the planet, its natural environment and humans' existence. It is clear the planet is warming, and warming rapidly. This is due to the greenhouse effect, which is caused by greenhouse gases, which trap heat inside the Earth's atmosphere because of their more complex molecular structure which allows them to vibrate and in turn trap heat and release it back towards the Earth. This warming is also responsible for the extinction of natural habitats, which in turn leads to a reduction in wildlife population. The most recent report from the Intergovernmental Panel on Climate Change (the group of the leading climate scientists in the world) concluded that the earth will warm anywhere from 2.7 to almost 11 degrees Fahrenheit (1.5 to 6 degrees Celsius) between 1990 and 2100. Efforts have been increasingly focused on the mitigation of greenhouse gases that are causing climatic changes, on developing adaptative strategies to global warming, to assist humans, other animal, and plant species, ecosystems, regions and nations in adjusting to the effects of global warming. Some examples of recent collaboration to address climate change and global warming include: The United Nations Framework Convention Treaty and convention on Climate Change, to stabilize greenhouse gas concentrations in the atmosphere at a level that would prevent dangerous anthropogenic interference with the climate system. The Kyoto Protocol, which is the protocol to the international Framework Convention on Climate Change treaty, again with the objective of reducing greenhouse gases in an effort to prevent anthropogenic climate change. The Western Climate Initiative, to identify, evaluate, and implement collective and cooperative ways to reduce greenhouse gases in the region, focusing on a market-based cap-and-trade system. A significantly profound challenge is to identify the natural environmental dynamics in contrast to environmental changes not within natural variances. A common solution is to adapt a static view neglecting natural variances to exist. Methodologically, this view could be defended when looking at processes which change slowly and short time series, while the problem arrives when fast processes turns essential in the object of the study. Climate Climate looks at the statistics of temperature, humidity, atmospheric pressure, wind, rainfall, atmospheric particle count and other meteorological elements in a given region over long periods of time. Weather, on the other hand, is the present condition of these same elements over periods up to two weeks. Climates can be classified according to the average and typical ranges of different variables, most commonly temperature and precipitation. The most commonly used classification scheme is the one originally developed by Wladimir Köppen. The Thornthwaite system, in use since 1948, uses evapotranspiration as well as temperature and precipitation information to study animal species diversity and the potential impacts of climate changes. Weather Weather is a set of all the phenomena occurring in a given atmospheric area at a given time. Most weather phenomena occur in the troposphere, just below the stratosphere. Weather refers, generally, to day-to-day temperature and precipitation activity, whereas climate is the term for the average atmospheric conditions over longer periods of time. When used without qualification, "weather" is understood to be the weather of Earth. Weather occurs due to density (temperature and moisture) differences between one place and another. These differences can occur due to the sun angle at any particular spot, which varies by latitude from the tropics. The strong temperature contrast between polar and tropical air gives rise to the jet stream. Weather systems in the mid-latitudes, such as extratropical cyclones, are caused by instabilities of the jet stream flow. Because the Earth's axis is tilted relative to its orbital plane, sunlight is incident at different angles at different times of the year. On the Earth's surface, temperatures usually range ±40 °C (100 °F to −40 °F) annually. Over thousands of years, changes in the Earth's orbit have affected the amount and distribution of solar energy received by the Earth and influenced long-term climate. Surface temperature differences in turn cause pressure differences. Higher altitudes are cooler than lower altitudes due to differences in compressional heating. Weather forecasting is the application of science and technology to predict the state of the atmosphere for a future time and a given location. The atmosphere is a chaotic system, and small changes to one part of the system can grow to have large effects on the system as a whole. Human attempts to control the weather have occurred throughout human history, and there is evidence that civilized human activity such as agriculture and industry has inadvertently modified weather patterns. Life Evidence suggests that life on Earth has existed for about 3.7 billion years. All known life forms share fundamental molecular mechanisms, and based on these observations, theories on the origin of life attempt to find a mechanism explaining the formation of a primordial single cell organism from which all life originates. There are many different hypotheses regarding the path that might have been taken from simple organic molecules via pre-cellular life to protocells and metabolism. Although there is no universal agreement on the definition of life, scientists generally accept that the biological manifestation of life is characterized by organization, metabolism, growth, adaptation, response to stimuli and reproduction. Life may also be said to be simply the characteristic state of organisms. In biology, the science of living organisms, "life" is the condition which distinguishes active organisms from inorganic matter, including the capacity for growth, functional activity and the continual change preceding death. A diverse variety of living organisms (life forms) can be found in the biosphere on Earth, and properties common to these organisms—plants, animals, fungi, protists, archaea, and bacteria—are a carbon- and water-based cellular form with complex organization and heritable genetic information. Living organisms undergo metabolism, maintain homeostasis, possess a capacity to grow, respond to stimuli, reproduce and, through natural selection, adapt to their environment in successive generations. More complex living organisms can communicate through various means. Ecosystems An ecosystem (also called an environment) is a natural unit consisting of all plants, animals, and micro-organisms (biotic factors) in an area functioning together with all of the non-living physical (abiotic) factors of the environment. Central to the ecosystem concept is the idea that living organisms are continually engaged in a highly interrelated set of relationships with every other element constituting the environment in which they exist. Eugene Odum, one of the founders of the science of ecology, stated: "Any unit that includes all of the organisms (i.e.: the "community") in a given area interacting with the physical environment so that a flow of energy leads to clearly defined trophic structure, biotic diversity, and material cycles (i.e.: exchange of materials between living and nonliving parts) within the system is an ecosystem." The human ecosystem concept is then grounded in the deconstruction of the human/nature dichotomy, and the emergent premise that all species are ecologically integrated with each other, as well as with the abiotic constituents of their biotope. A more significant number or variety of species or biological diversity of an ecosystem may contribute to greater resilience of an ecosystem because there are more species present at a location to respond to change and thus "absorb" or reduce its effects. This reduces the effect before the ecosystem's structure changes to a different state. This is not universally the case and there is no proven relationship between the species diversity of an ecosystem and its ability to provide goods and services on a sustainable level. The term ecosystem can also pertain to human-made environments, such as human ecosystems and human-influenced ecosystems. It can describe any situation where there is relationship between living organisms and their environment. Fewer areas on the surface of the earth today exist free from human contact, although some genuine wilderness areas continue to exist without any forms of human intervention. Biogeochemical cycles Global biogeochemical cycles are critical to life, most notably those of water, oxygen, carbon, nitrogen and phosphorus. The nitrogen cycle is the transformation of nitrogen and nitrogen-containing compounds in nature. It is a cycle which includes gaseous components. The water cycle, is the continuous movement of water on, above, and below the surface of the Earth. Water can change states among liquid, vapour, and ice at various places in the water cycle. Although the balance of water on Earth remains fairly constant over time, individual water molecules can come and go. The carbon cycle is the biogeochemical cycle by which carbon is exchanged among the biosphere, pedosphere, geosphere, hydrosphere, and atmosphere of the Earth. The oxygen cycle is the movement of oxygen within and between its three main reservoirs: the atmosphere, the biosphere, and the lithosphere. The main driving factor of the oxygen cycle is photosynthesis, which is responsible for the modern Earth's atmospheric composition and life. The phosphorus cycle is the movement of phosphorus through the lithosphere, hydrosphere, and biosphere. The atmosphere does not play a significant role in the movements of phosphorus, because phosphorus and phosphorus compounds are usually solids at the typical ranges of temperature and pressure found on Earth. Wilderness Wilderness is generally defined as a natural environment on Earth that has not been significantly modified by human activity. The WILD Foundation goes into more detail, defining wilderness as: "The most intact, undisturbed wild natural areas left on our planet – those last truly wild places that humans do not control and have not developed with roads, pipelines or other industrial infrastructure." Wilderness areas and protected parks are considered important for the survival of certain species, ecological studies, conservation, solitude, and recreation. Wilderness is deeply valued for cultural, spiritual, moral, and aesthetic reasons. Some nature writers believe wilderness areas are vital for the human spirit and creativity. The word, "wilderness", derives from the notion of wildness; in other words that which is not controllable by humans. The word etymology is from the Old English wildeornes, which in turn derives from wildeor meaning wild beast (wild + deor = beast, deer). From this point of view, it is the wildness of a place that makes it a wilderness. The mere presence or activity of people does not disqualify an area from being "wilderness". Many ecosystems that are, or have been, inhabited or influenced by activities of people may still be considered "wild". This way of looking at wilderness includes areas within which natural processes operate without very noticeable human interference. Wildlife includes all non-domesticated plants, animals and other organisms. Domesticating wild plant and animal species for human benefit has occurred many times all over the planet, and has a major impact on the environment, both positive and negative. Wildlife can be found in all ecosystems. Deserts, rain forests, plains, and other areas—including the most developed urban sites—all have distinct forms of wildlife. While the term in popular culture usually refers to animals that are untouched by civilized human factors, most scientists agree that wildlife around the world is (now) impacted by human activities. Challenges It is the common understanding of natural environment that underlies environmentalism — a broad political, social and philosophical movement that advocates various actions and policies in the interest of protecting what nature remains in the natural environment, or restoring or expanding the role of nature in this environment. While true wilderness is increasingly rare, wild nature (e.g., unmanaged forests, uncultivated grasslands, wildlife, wildflowers) can be found in many locations previously inhabited by humans. Goals for the benefit of people and natural systems, commonly expressed by environmental scientists and environmentalists include: Elimination of pollution and toxicants in air, water, soil, buildings, manufactured goods, and food. Preservation of biodiversity and protection of endangered species. Conservation and sustainable use of resources such as water, land, air, energy, raw materials, and natural resources. Halting human-induced global warming, which represents pollution, a threat to biodiversity, and a threat to human populations. Shifting from fossil fuels to renewable energy in electricity, heating and cooling, and transportation, which addresses pollution, global warming, and sustainability. This may include public transportation and distributed generation, which have benefits for traffic congestion and electric reliability. Shifting from meat-intensive diets to largely plant-based diets in order to help mitigate biodiversity loss and climate change. Establishment of nature reserves for recreational purposes and ecosystem preservation. Sustainable and less polluting waste management including waste reduction (or even zero waste), reuse, recycling, composting, waste-to-energy, and anaerobic digestion of sewage sludge. Reducing profligate consumption and clamping down on illegal fishing and logging. Slowing and stabilisation of human population growth. Reducing the import of second hand electronic appliances from developed countries to developing countries. Criticism In some cultures the term environment is meaningless because there is no separation between people and what they view as the natural world, or their surroundings. Specifically in the United States and Arabian countries many native cultures do not recognize the "environment", or see themselves as environmentalists.	Physical sciences	Earth science basics: General	Earth science
558959	https://en.wikipedia.org/wiki/Railway%20electrification	Railway electrification	Railway electrification is the use of electric power for the propulsion of rail transport. Electric railways use either electric locomotives (hauling passengers or freight in separate cars), electric multiple units (passenger cars with their own motors) or both. Electricity is typically generated in large and relatively efficient generating stations, transmitted to the railway network and distributed to the trains. Some electric railways have their own dedicated generating stations and transmission lines, but most purchase power from an electric utility. The railway usually provides its own distribution lines, switches, and transformers. Power is supplied to moving trains with a (nearly) continuous conductor running along the track that usually takes one of two forms: an overhead line, suspended from poles or towers along the track or from structure or tunnel ceilings and contacted by a pantograph, or a third rail mounted at track level and contacted by a sliding "pickup shoe". Both overhead wire and third-rail systems usually use the running rails as the return conductor, but some systems use a separate fourth rail for this purpose. In comparison to the principal alternative, the diesel engine, electric railways offer substantially better energy efficiency, lower emissions, and lower operating costs. Electric locomotives are also usually quieter, more powerful, and more responsive and reliable than diesel. They have no local emissions, an important advantage in tunnels and urban areas. Some electric traction systems provide regenerative braking that turns the train's kinetic energy back into electricity and returns it to the supply system to be used by other trains or the general utility grid. While diesel locomotives burn petroleum products, electricity can be generated from diverse sources, including renewable energy. Historically, concerns of resource independence have played a role in the decision to electrify railway lines. The landlocked Swiss confederation which almost completely lacks oil or coal deposits but has plentiful hydropower electrified its network in part in reaction to supply issues during both World Wars. Disadvantages of electric traction include: high capital costs that may be uneconomic on lightly trafficked routes, a relative lack of flexibility (since electric trains need third rails or overhead wires), and a vulnerability to power interruptions. Electro-diesel locomotives and electro-diesel multiple units mitigate these problems somewhat as they are capable of running on diesel power during an outage or on non-electrified routes. Different regions may use different supply voltages and frequencies, complicating through service and requiring greater complexity of locomotive power. There used to be a historical concern for double-stack rail transport regarding clearances with overhead lines but it is no longer universally true , with both Indian Railways and China Railway regularly operating electric double-stack cargo trains under overhead lines. Railway electrification has constantly increased in the past decades, and as of 2022, electrified tracks account for nearly one-third of total tracks globally. History Railway electrification is the development of powering trains and locomotives using electricity instead of diesel or steam power. The history of railway electrification dates back to the late 19th century when the first electric tramways were introduced in cities like Berlin, London, and New York City. In 1881, the first permanent railway electrification in the world was the Gross-Lichterfelde Tramway in Berlin, Germany. Overhead line electrification was first applied successfully by Frank Sprague in Richmond, Virginia in 1887-1888, and led to the electrification of hundreds of additional street railway systems by the early 1890s. The first electrification of a mainline railway was the Baltimore and Ohio Railroad's Baltimore Belt Line in the United States in 1895–96. The early electrification of railways used direct current (DC) power systems, which were limited in terms of the distance they could transmit power. However, in the early 20th century, alternating current (AC) power systems were developed, which allowed for more efficient power transmission over longer distances. In the 1920s and 1930s, many countries worldwide began to electrify their railways. In Europe, Switzerland, Sweden, France, and Italy were among the early adopters of railway electrification. In the United States, the New York, New Haven and Hartford Railroad was one of the first major railways to be electrified. Railway electrification continued to expand throughout the 20th century, with technological improvements and the development of high-speed trains and commuters. Today, many countries have extensive electrified railway networks with of standard lines in the world, including China, India, Japan, France, Germany, and the United Kingdom. Electrification is seen as a more sustainable and environmentally friendly alternative to diesel or steam power and is an important part of many countries' transportation infrastructure. Classification Electrification systems are classified by three main parameters: Voltage Current Direct current (DC) Alternating current (AC) Frequency Contact system Overhead lines (catenary) Third rail Fourth rail Ground-level power supply Selection of an electrification system is based on economics of energy supply, maintenance, and capital cost compared to the revenue obtained for freight and passenger traffic. Different systems are used for urban and intercity areas; some electric locomotives can switch to different supply voltages to allow flexibility in operation. Standardised voltages Six of the most commonly used voltages have been selected for European and international standardisation. Some of these are independent of the contact system used, so that, for example, 750VDC may be used with either third rail or overhead lines. There are many other voltage systems used for railway electrification systems around the world, and the list of railway electrification systems covers both standard voltage and non-standard voltage systems. The permissible range of voltages allowed for the standardised voltages is as stated in standards BSEN50163 and IEC60850. These take into account the number of trains drawing current and their distance from the substation. Direct current Overhead lines 1,500V DC is used in Japan, Indonesia, Hong Kong (parts), Ireland, Australia (parts), France (also using , the Netherlands, New Zealand (Wellington), Singapore (on the North East MRT line), the United States (Chicago area on the Metra Electric district and the South Shore Line interurban line and Link light rail in Seattle, Washington). In Slovakia, there are two narrow-gauge lines in the High Tatras (one a cog railway). In the Netherlands it is used on the main system, alongside 25kV on the HSL-Zuid and Betuwelijn, and 3,000V south of Maastricht. In Portugal, it is used in the Cascais Line and in Denmark on the suburban S-train system (1650V DC). In the United Kingdom, 1,500VDC was used in 1954 for the Woodhead trans-Pennine route (now closed); the system used regenerative braking, allowing for transfer of energy between climbing and descending trains on the steep approaches to the tunnel. The system was also used for suburban electrification in East London and Manchester, now converted to 25kVAC. It is now only used for the Tyne and Wear Metro. In India, 1,500V DC was the first electrification system launched in 1925 in Mumbai area. Between 2012 and 2016, the electrification was converted to 25kV 50Hz, which is the countrywide system. 3kV DC is used in Belgium, Italy, Spain, Poland, Slovakia, Slovenia, South Africa, Chile, the northern portion of the Czech Republic, the former republics of the Soviet Union, and in the Netherlands on a few kilometers between Maastricht and Belgium. It was formerly used by the Milwaukee Road from Harlowton, Montana, to Seattle, across the Continental Divide and including extensive branch and loop lines in Montana, and by the Delaware, Lackawanna and Western Railroad (now New Jersey Transit, converted to 25kVAC) in the United States, and the Kolkata suburban railway (Bardhaman Main Line) in India, before it was converted to 25kV 50Hz. DC voltages between 600V and 750V are used by most tramways and trolleybus networks, as well as some metro systems as the traction motors accept this voltage without the weight of an on-board transformer. Medium-voltage DC Increasing availability of high-voltage semiconductors may allow the use of higher and more efficient DC voltages that heretofore have only been practical with AC. The use of medium-voltage DC electrification (MVDC) would solve some of the issues associated with standard-frequency AC electrification systems, especially possible supply grid load imbalance and the phase separation between the electrified sections powered from different phases, whereas high voltage would make the transmission more efficient. UIC conducted a case study for the conversion of the Bordeaux-Hendaye railway line (France), currently electrified at 1.5kV DC, to 9kV DC and found that the conversion would allow to use less bulky overhead wires (saving €20 million per 100route-km) and lower the losses (saving 2GWh per year per 100route-km; equalling about €150,000 p.a.). The line chosen is one of the lines, totalling 6000km, that are in need of renewal. In the 1960s the Soviets experimented with boosting the overhead voltage from 3 to 6kV. DC rolling stock was equipped with ignitron-based converters to lower the supply voltage to 3kV. The converters turned out to be unreliable and the experiment was curtailed. In 1970 the Ural Electromechanical Institute of Railway Engineers carried out calculations for railway electrification at , showing that the equivalent loss levels for a system could be achieved with DC voltage between 11 and 16kV. In the 1980s and 1990s was being tested on the October Railway near Leningrad (now Petersburg). The experiments ended in 1995 due to the end of funding. Third rail Most electrification systems use overhead wires, but third rail is an option up to 1,500V. Third rail systems almost exclusively use DC distribution. The use of AC is usually not feasible due to the dimensions of a third rail being physically very large compared with the skin depth that AC penetrates to in a steel rail. This effect makes the resistance per unit length unacceptably high compared with the use of DC. Third rail is more compact than overhead wires and can be used in smaller-diameter tunnels, an important factor for subway systems. Fourth rail The London Underground in England is one of few networks that uses a four-rail system. The additional rail carries the electrical return that, on third-rail and overhead networks, is provided by the running rails. On the London Underground, a top-contact third rail is beside the track, energized at , and a top-contact fourth rail is located centrally between the running rails at , which combine to provide a traction voltage of . The same system was used for Milan's earliest underground line, Milan Metro's line 1, whose more recent lines use an overhead catenary or a third rail. The key advantage of the four-rail system is that neither running rail carries any current. This scheme was introduced because of the problems of return currents, intended to be carried by the earthed (grounded) running rail, flowing through the iron tunnel linings instead. This can cause electrolytic damage and even arcing if the tunnel segments are not electrically bonded together. The problem was exacerbated because the return current also had a tendency to flow through nearby iron pipes forming the water and gas mains. Some of these, particularly Victorian mains that predated London's underground railways, were not constructed to carry currents and had no adequate electrical bonding between pipe segments. The four-rail system solves the problem. Although the supply has an artificially created earth point, this connection is derived by using resistors which ensures that stray earth currents are kept to manageable levels. Power-only rails can be mounted on strongly insulating ceramic chairs to minimise current leak, but this is not possible for running rails, which have to be seated on stronger metal chairs to carry the weight of trains. However, elastomeric rubber pads placed between the rails and chairs can now solve part of the problem by insulating the running rails from the current return should there be a leakage through the running rails. The Expo and Millennium Line of the Vancouver SkyTrain use side-contact fourth-rail systems for their supply. Both are located to the side of the train, as the space between the running rails is occupied by an aluminum plate, as part of stator of the linear induction propulsion system used on the Innovia ART system. While part of the SkyTrain network, the Canada Line does not use this system and instead uses more traditional motors attached to the wheels and third-rail electrification. Rubber-tyred systems A few lines of the Paris Métro in France operate on a four-rail power system. The trains move on rubber tyres which roll on a pair of narrow roll ways made of steel and, in some places, of concrete. Since the tyres do not conduct the return current, the two guide bars provided outside the running 'roll ways' become, in a sense, a third and fourth rail which each provide , so at least electrically it is a four-rail system. Each wheel set of a powered bogie carries one traction motor. A side sliding (side running) contact shoe picks up the current from the vertical face of each guide bar. The return of each traction motor, as well as each wagon, is effected by one contact shoe each that slide on top of each one of the running rails. This and all other rubber-tyred metros that have a track between the roll ways operate in the same manner. Alternating current Railways and electrical utilities use AC as opposed to DC for the same reason: to use transformers, which require AC, to produce higher voltages. The higher the voltage, the lower the current for the same power (because power is current multiplied by voltage), and power loss is proportional to the current squared. The lower current reduces line loss, thus allowing higher power to be delivered. As alternating current is used with high voltages. Inside the locomotive, a transformer steps the voltage down for use by the traction motors and auxiliary loads. An early advantage of AC is that the power-wasting resistors used in DC locomotives for speed control were not needed in an AC locomotive: multiple taps on the transformer can supply a range of voltages. Separate low-voltage transformer windings supply lighting and the motors driving auxiliary machinery. More recently, the development of very high power semiconductors has caused the classic DC motor to be largely replaced with the three-phase induction motor fed by a variable frequency drive, a special inverter that varies both frequency and voltage to control motor speed. These drives can run equally well on DC or AC of any frequency, and many modern electric locomotives are designed to handle different supply voltages and frequencies to simplify cross-border operation. Low-frequency alternating current Five European countries Germany, Austria, Switzerland, Norway and Sweden have standardized on 15kV Hz (the 50Hz mains frequency divided by three) single-phase AC. On 16 October 1995, Germany, Austria and Switzerland changed from Hz to 16.7Hz which is no longer exactly one-third of the grid frequency. This solved overheating problems with the rotary converters used to generate some of this power from the grid supply. In the US, the New York, New Haven, and Hartford Railroad, the Pennsylvania Railroad and the Philadelphia and Reading Railway adopted 11kV 25Hz single-phase AC. Parts of the original electrified network still operate at 25Hz, with voltage boosted to 12kV, while others were converted to 12.5 or 25kV 60Hz. In the UK, the London, Brighton and South Coast Railway pioneered overhead electrification of its suburban lines in London, London Bridge to Victoria being opened to traffic on 1December 1909. Victoria to Crystal Palace via Balham and West Norwood opened in May 1911. Peckham Rye to West Norwood opened in June 1912. Further extensions were not made owing to the First World War. Two lines opened in 1925 under the Southern Railway serving Coulsdon North and Sutton railway station. The lines were electrified at 6.7kV 25Hz. It was announced in 1926 that all lines were to be converted to DC third rail and the last overhead-powered electric service ran in September 1929. Standard frequency alternating current AC power is used at 60Hz in North America (excluding the aforementioned 25Hz network), western Japan, South Korea and Taiwan; and at 50Hz in a number of European countries, India, Saudi Arabia, eastern Japan, countries that used to be part of the Soviet Union, on high-speed lines in much of Western Europe (including countries that still run conventional railways under DC but not in countries using 16.7Hz, see above). Most systems like this operate at 25kV, although 12.5kV sections exist in the United States, and 20kV is used on some narrow-gauge lines in Japan. On "French system" HSLs, the overhead line and a "sleeper" feeder line each carry 25kV in relation to the rails, but in opposite phase so they are at 50kV from each other; autotransformers equalize the tension at regular intervals. Three-phase alternating current Various railway electrification systems in the late nineteenth and twentieth centuries utilised three-phase, rather than single-phase electric power delivery due to ease of design of both power supply and locomotives. These systems could either use standard network frequency and three power cables, or reduced frequency, which allowed for return-phase line to be third rail, rather than an additional overhead wire. Comparisons AC versus DC for mainlines The majority of modern electrification systems take AC energy from a power grid that is delivered to a locomotive, and within the locomotive, transformed and rectified to a lower DC voltage in preparation for use by traction motors. These motors may either be DC motors which directly use the DC or they may be three-phase AC motors which require further conversion of the DC to variable frequency three-phase AC (using power electronics). Thus both systems are faced with the same task: converting and transporting high-voltage AC from the power grid to low-voltage DC in the locomotive. The difference between AC and DC electrification systems lies in where the AC is converted to DC: at the substation or on the train. Energy efficiency and infrastructure costs determine which of these is used on a network, although this is often fixed due to pre-existing electrification systems. Both the transmission and conversion of electric energy involve losses: ohmic losses in wires and power electronics, magnetic field losses in transformers and smoothing reactors (inductors). Power conversion for a DC system takes place mainly in a railway substation where large, heavy, and more efficient hardware can be used as compared to an AC system where conversion takes place aboard the locomotive where space is limited and losses are significantly higher. However, the higher voltages used in many AC electrification systems reduce transmission losses over longer distances, allowing for fewer substations or more powerful locomotives to be used. Also, the energy used to blow air to cool transformers, power electronics (including rectifiers), and other conversion hardware must be accounted for. Standard AC electrification systems use much higher voltages than standard DC systems. One of the advantages of raising the voltage is that, to transmit certain level of power, lower current is necessary (). Lowering the current reduces the ohmic losses and allows for less bulky, lighter overhead line equipment and more spacing between traction substations, while maintaining power capacity of the system. On the other hand, the higher voltage requires larger isolation gaps, requiring some elements of infrastructure to be larger. The standard-frequency AC system may introduce imbalance to the supply grid, requiring careful planning and design (as at each substation power is drawn from two out of three phases). The low-frequency AC system may be powered by separate generation and distribution network or a network of converter substations, adding the expense, also low-frequency transformers, used both at the substations and on the rolling stock, are particularly bulky and heavy. The DC system, apart from being limited as to the maximum power that can be transmitted, also can be responsible for electrochemical corrosion due to stray DC currents. Electric versus diesel Energy efficiency Electric trains need not carry the weight of prime movers, transmission and fuel. This is partly offset by the weight of electrical equipment. Regenerative braking returns power to the electrification system so that it may be used elsewhere, by other trains on the same system or returned to the general power grid. This is especially useful in mountainous areas where heavily loaded trains must descend long grades. Central station electricity can often be generated with higher efficiency than a mobile engine/generator. While the efficiency of power plant generation and diesel locomotive generation are roughly the same in the nominal regime, diesel motors decrease in efficiency in non-nominal regimes at low power while if an electric power plant needs to generate less power it will shut down its least efficient generators, thereby increasing efficiency. The electric train can save energy (as compared to diesel) by regenerative braking and by not needing to consume energy by idling as diesel locomotives do when stopped or coasting. However, electric rolling stock may run cooling blowers when stopped or coasting, thus consuming energy. Large fossil fuel power stations operate at high efficiency, and can be used for district heating or to produce district cooling, leading to a higher total efficiency. Electricity for electric rail systems can also come from renewable energy, nuclear power, or other low-carbon sources, which do not emit pollution or emissions. Power output Electric locomotives may easily be constructed with greater power output than most diesel locomotives. For passenger operation it is possible to provide enough power with diesel engines (see e.g. 'ICE TD') but, at higher speeds, this proves costly and impractical. Therefore, almost all high speed trains are electric. The high power of electric locomotives also gives them the ability to pull freight at higher speed over gradients; in mixed traffic conditions this increases capacity when the time between trains can be decreased. The higher power of electric locomotives and an electrification can also be a cheaper alternative to a new and less steep railway if train weights are to be increased on a system. On the other hand, electrification may not be suitable for lines with low frequency of traffic, because lower running cost of trains may be outweighed by the high cost of the electrification infrastructure. Therefore, most long-distance lines in developing or sparsely populated countries are not electrified due to relatively low frequency of trains. Network effect Network effects are a large factor with electrification. When converting lines to electric, the connections with other lines must be considered. Some electrifications have subsequently been removed because of the through traffic to non-electrified lines. If through traffic is to have any benefit, time-consuming engine switches must occur to make such connections or expensive dual mode engines must be used. This is mostly an issue for long-distance trips, but many lines come to be dominated by through traffic from long-haul freight trains (usually running coal, ore, or containers to or from ports). In theory, these trains could enjoy dramatic savings through electrification, but it can be too costly to extend electrification to isolated areas, and unless an entire network is electrified, companies often find that they need to continue use of diesel trains even if sections are electrified. The increasing demand for container traffic, which is more efficient when utilizing the double-stack car, also has network effect issues with existing electrifications due to insufficient clearance of overhead electrical lines for these trains, but electrification can be built or modified to have sufficient clearance, at additional cost. A problem specifically related to electrified lines are gaps in the electrification. Electric vehicles, especially locomotives, lose power when traversing gaps in the supply, such as phase change gaps in overhead systems, and gaps over points in third rail systems. These become a nuisance if the locomotive stops with its collector on a dead gap, in which case there is no power to restart. This is less of a problem in trains consisting of two or more multiple units coupled together, since in that case if the train stops with one collector in a dead gap, another multiple unit can push or pull the disconnected unit until it can again draw power. The same applies to the kind of push-pull trains which have a locomotive at each end. Power gaps can be overcome in single-collector trains by on-board batteries or motor-flywheel-generator systems. In 2014, progress is being made in the use of large capacitors to power electric vehicles between stations, and so avoid the need for overhead wires between those stations. Maintenance costs Maintenance costs of the lines may be increased by electrification, but many systems claim lower costs due to reduced wear-and-tear on the track from lighter rolling stock. There are some additional maintenance costs associated with the electrical equipment around the track, such as power sub-stations and the catenary wire itself, but, if there is sufficient traffic, the reduced track and especially the lower engine maintenance and running costs exceed the costs of this maintenance significantly. Sparks effect Newly electrified lines often show a "sparks effect", whereby electrification in passenger rail systems leads to significant jumps in patronage / revenue. The reasons may include electric trains being seen as more modern and attractive to ride, faster, quieter and smoother service, and the fact that electrification often goes hand in hand with a general infrastructure and rolling stock overhaul / replacement, which leads to better service quality (in a way that theoretically could also be achieved by doing similar upgrades yet without electrification). Whatever the causes of the sparks effect, it is well established for numerous routes that have electrified over decades. This also applies when bus routes with diesel buses are replaced by trolleybuses. The overhead wires make the service "visible" even in no bus is running and the existence of the infrastructure gives some long-term expectations of the line being in operation. Double-stack rail transport Due to the height restriction imposed by the overhead wires, double-stacked container trains have been traditionally difficult and rare to operate under electrified lines. However, this limitation is being overcome by railways in India, China and African countries by laying new tracks with increased catenary height. Such installations are in the Western Dedicated Freight Corridor in India where the wire height is at to accommodate double-stack container trains without the need of well-wagons. Advantages There are a number of advantages including the fact there is no exposure of passengers to exhaust from the locomotive and lower cost of building, running and maintaining locomotives and multiple units. Electric trains have a higher power-to-weight ratio (no onboard fuel tanks), resulting in fewer locomotives, faster acceleration, higher practical limit of power, higher limit of speed, less noise pollution (quieter operation). The faster acceleration clears lines more quickly to run more trains on the track in urban rail uses. Reduced power loss at higher altitudes (for power loss see Diesel engine) Independence of running costs from fluctuating fuel prices Service to underground stations where diesel trains cannot operate for safety reasons Reduced environmental pollution, especially in highly populated urban areas, even if electricity is produced by fossil fuels Easily accommodates kinetic energy brake reclaim using supercapacitors More comfortable ride on multiple units as trains have no underfloor diesel engines Somewhat higher energy efficiency in part due to regenerative braking and less power lost when "idling" More flexible primary energy source: can use coal, natural gas, nuclear or renewable energy (hydro, solar, wind) as the primary energy source instead of diesel fuel If the entire network is electrified, diesel infrastructure such as fueling stations, maintenance yards and indeed the diesel locomotive fleet can be retired or put to other uses – this is often the business case in favor of electrifying the last few lines in a network where otherwise costs would be too high. Having only one type of motive power also allows greater fleet homogeneity which can also reduce costs. Disadvantages Electrification cost: electrification requires an entire new infrastructure to be built around the existing tracks at a significant cost. Costs are especially high when tunnels, bridges and other obstructions have to be altered for clearance. Another aspect that can raise the cost of electrification are the alterations or upgrades to railway signalling needed for new traffic characteristics, and to protect signalling circuitry and track circuits from interference by traction current. Electrification typically requires line closures while new equipment is being installed. Appearance: the overhead line structures and cabling can have a significant landscape impact compared with a non-electrified or third rail electrified line that has only occasional signalling equipment above ground level. Fragility and vulnerability: overhead electrification systems can suffer severe disruption due to minor mechanical faults or the effects of high winds causing the pantograph of a moving train to become entangled with the catenary, ripping the wires from their supports. The damage is often not limited to the supply to one track, but extends to those for adjacent tracks as well, causing the entire route to be blocked for a considerable time. Third-rail systems can suffer disruption in cold weather due to ice forming on the conductor rail. Theft: the high scrap value of copper and the unguarded, remote installations make overhead cables an attractive target for scrap metal thieves. Attempts at theft of live 25kV cables may end in the thief's death from electrocution. In the UK, cable theft is claimed to be one of the biggest sources of delay and disruption to train services – though this normally relates to signalling cable, which is equally problematic for diesel lines. Incompatibility: Diesel trains can run on any track without electricity or with any kind of electricity (third rail or overhead line, DC or AC, and at any voltage or frequency). Not so for electric trains, which can never run on non-electrified lines, and which even on electrified lines can run only on the single, or the few, electrical system(s) for which they are equipped. Even on fully electrified networks, it is usually a good idea to keep a few diesel locomotives for maintenance and repair trains, for instance to repair broken or stolen overhead lines, or to lay new tracks. However, due to ventilation issues, diesel trains may have to be banned from certain tunnels and underground train stations mitigating the advantage of diesel trains somewhat. Birds may perch on parts with different charges, and animals may also touch the electrification system. Dead animals attract foxes or other scavengers, bringing risk of collision with trains. In most of the world's railway networks, the height clearance of overhead electrical lines is not sufficient for a double-stack container car or other unusually tall loads. To upgrade electrified lines to the correct clearances () to take double-stacked container trains, besides renewing bridges over it, would normally mean need for special pantographs violating standardisation and requiring custom made vehicles. Railway electrification around the world As of 2012, electrified tracks accounted for nearly one third of total tracks globally. As of 2018, there were of railways electrified at 25kV, either 50 or 60Hz; electrified at ; electrified at 15kV 16.7 or Hz and electrified at . As of 2023, the Swiss rail network is the largest fully electrified network in the world and one of only eleven countries or territories to achieve this, as listed in List of countries by rail transport network size. The percentage then continues falling in order with Laos, Montenegro, India, Belgium, Georgia, South Korea, Netherlands, and Japan, with all others being less than 75% electrified. Overall, China takes first place, with around of electrified railway, followed by India with over of electrified railway, and continuing with Russia, with over of electrified railway. A number of countries have zero electrified railways, instead relying on diesel multiple units, locomotive hauled services and many alternate forms of transport. The European Union contains the longest amount of electrified railways (in length), with over of electrified railway, however only making up around 55% of the total railway length. Several countries have announced plans to electrify all or most of their railway network, including Indian Railways and Israel Railways. The Trans-Siberian Railway mainly in Russia is completely electrified, making it one of the longest stretches of electrified railways in the world.	Technology	Rail and cable transport	null
559011	https://en.wikipedia.org/wiki/Lactococcus%20lactis	Lactococcus lactis	Lactococcus lactis is a gram-positive bacterium used extensively in the production of buttermilk and cheese, but has also become famous as the first genetically modified organism to be used alive for the treatment of human disease. L. lactis cells are cocci that group in pairs and short chains, and, depending on growth conditions, appear ovoid with a typical length of 0.5 - 1.5 μm. L. lactis does not produce spores (nonsporulating) and are not motile (nonmotile). They have a homofermentative metabolism, meaning they produce lactic acid from sugars. They've also been reported to produce exclusive L-(+)-lactic acid. However, reported D-(−)-lactic acid can be produced when cultured at low pH. The capability to produce lactic acid is one of the reasons why L. lactis is one of the most important microorganisms in the dairy industry. Based on its history in food fermentation, L. lactis has generally recognized as safe (GRAS) status, with few case reports of it being an opportunistic pathogen. Lactococcus lactis is of crucial importance for manufacturing dairy products, such as buttermilk and cheeses. When L. lactis ssp. lactis is added to milk, the bacterium uses enzymes to produce energy molecules (ATP), from lactose. The byproduct of ATP energy production is lactic acid. The lactic acid produced by the bacterium curdles the milk, which then separates to form curds that are used to produce cheese. Other uses that have been reported for this bacterium include the production of pickled vegetables, beer or wine, some breads, and other fermented foodstuffs like soymilk kefir, buttermilk, and others. L. lactis is one of the best characterized low GC Gram positive bacteria with detailed knowledge on genetics, metabolism and biodiversity. L. lactis is mainly isolated from either the dairy environment, or plant material. Dairy isolates are suggested to have evolved from plant isolates through a process in which genes without benefit in the rich milk were lost or downregulated. This process, called genome erosion or reductive evolution, has been described in several other lactic acid bacteria. The proposed transition from the plant to the dairy environment was reproduced in the laboratory through experimental evolution of a plant isolate that was cultivated in milk for a prolonged period. Consistent with the results from comparative genomics (see references above), this resulted in L. lactis losing or downregulating genes that are dispensable in milk and the upregulation of peptide transport. Hundreds of novel small RNAs were identified by Meulen et al. in the genome of L. lactis MG1363. One of them, LLnc147, was shown to be involved in carbon uptake and metabolism. Cheese production L. lactis subsp. lactis (formerly Streptococcus lactis) is used in the early stages for the production of many cheeses, including brie, camembert, Cheddar, Colby, Gruyère, Parmesan, and Roquefort. The use of L. lactis in dairy factories is not without issues. Bacteriophages specific to L. lactis cause significant economic losses each year by preventing the bacteria from fully metabolizing the milk substrate. Several epidemiologic studies showed the phages mainly responsible for these losses are from the species 936, c2, and P335 (all from the family Siphoviridae). The state Assembly of Wisconsin, also the number one cheese-producing state in the United States, voted in 2010 to name this bacterium as the official state microbe; it would have been the first and only such designation by a state legislature in the nation, however the legislation was not adopted by the Senate. The legislation was introduced in November 2009 as Assembly Bill 556 by Representatives Hebl, Vruwink, Williams, Pasch, Danou, and Fields; it was cosponsored by Senator Taylor. The bill passed the Assembly on May 15, 2010, and was dropped by the Senate on April 28. Therapeutic benefits The feasibility of using lactic acid bacteria (LAB) as functional protein delivery vectors has been widely investigated. Lactococcus lactis has been demonstrated to be a promising candidate for the delivery of functional proteins because of its noninvasive and nonpathogenic characteristics. Many different expression systems of L. lactis have been developed and used for heterologous protein expression. Lactose fermentation In one study that sought to prove that some fermentation produced by L. lactis can hinder motility in pathogenic bacteria, the motilities of Pseudomonas, Vibrio, and Leptospira strains were severely disrupted by lactose utilization on the part of L. lactis. Using flagellar Salmonella as the experimental group, the research team found that a product of lactose fermentation is the cause of motility impairment in Salmonella. It is suggested that the L. lactis supernatant mainly affects Salmonella motility through disruption of flagellar rotation rather than through irreversible damage to morphology and physiology. Lactose fermentation by L. lactis produces acetate that reduces the intracellular pH of Salmonella, which in turn slows the rotation of their flagella. These results highlight the potential use of L. lactis for preventing infections by multiple bacterial species. Secretion of Interleukin-10 Genetically engineered L. lactis can secrete the cytokine interleukin-10 (IL-10) for the treatment of inflammatory bowel diseases (IBD), since IL-10 has a central role in downregulating inflammatory cascades and matrix metalloproteinases. A study by Lothar Steidler and Wolfgang Hans shows that in situ synthesis of IL-10 by genetically engineered L. lactis requires much lower doses than systemic treatments like antibodies to tumor necrosis factor (TNF) or recombinant IL-10. The authors propose two possible routes by which IL-10 can reach its therapeutic target. Genetically engineered L. lactis may produce murine IL-10 in the lumen, and the protein may diffuse to responsive cells in the epithelium or the lamina propria. Another route involves L. lactis being taken up by M cells because of its bacterial size and shape, and the major part of the effect may be due to recombinant IL-10 production in situ in intestinal lymphoid tissue. Both routes may involve paracellular transport mechanisms that are enhanced in inflammation. After transport, IL-10 may directly downregulate inflammation. In principle, this method may be useful for intestinal delivery of other protein therapeutics that are unstable or difficult to produce in large quantities and an alternative to the systemic treatment of IBD. Tumor-suppressor through Tumor metastasis-inhibiting peptide KISS1 Another study, led by Zhang B, created a L. lactis strain that maintains a plasmid containing a tumor metastasis-inhibiting peptide known as KISS1. L. lactis NZ9000 was demonstrated to be a cell factory for the secretion of biologically active KiSS1 protein, exerting inhibition effects on human colorectal cancer HT-29 cells. KiSS1 secreted from recombinant L. lactis strain effectively downregulated the expression of Matrix metalloproteinases (MMP-9), a crucial key in the invasion, metastasis, and regulation of the signaling pathways controlling tumor cell growth, survival, invasion, inflammation, and angiogenesis. The reason for this is that KiSS1 expressed in L. lactis activates the MAPK pathway via GPR54 signaling, suppressing NFκB binding to the MMP-9 promoter and thus downregulating MMP-9 expression. This, in turn, reduces the survival rate, inhibits metastasis, and induces dormancy of cancer cells. In addition, it was demonstrated that tumor growth can be inhibited by the LAB strain itself, due to the ability of LAB to produce exopolysaccharides. This study shows that L. lactis NZ9000 can inhibit HT-29 proliferation and induce cell apoptosis by itself. The success of this strain's construction helped to inhibit migration and expansion of cancer cells, showing that the secretion properties of L. lactis of this particular peptide may serve as a new tool for cancer therapy in the future.	Biology and health sciences	Gram-positive bacteria	Plants
560502	https://en.wikipedia.org/wiki/Chromosomal%20translocation	Chromosomal translocation	In genetics, chromosome translocation is a phenomenon that results in unusual rearrangement of chromosomes. This includes balanced and unbalanced translocation, with two main types: reciprocal, and Robertsonian translocation. Reciprocal translocation is a chromosome abnormality caused by exchange of parts between non-homologous chromosomes. Two detached fragments of two different chromosomes are switched. Robertsonian translocation occurs when two non-homologous chromosomes get attached, meaning that given two healthy pairs of chromosomes, one of each pair "sticks" and blends together homogeneously. A gene fusion may be created when the translocation joins two otherwise-separated genes. It is detected on cytogenetics or a karyotype of affected cells. Translocations can be balanced (in an even exchange of material with no genetic information extra or missing, and ideally full functionality) or unbalanced (where the exchange of chromosome material is unequal resulting in extra or missing genes). Reciprocal translocations Reciprocal translocations are usually an exchange of material between non-homologous chromosomes and occur in about 1 in 491 live births. Such translocations are usually harmless, as they do not result in a gain or loss of genetic material, though they may be detected in prenatal diagnosis. However, carriers of balanced reciprocal translocations may create gametes with unbalanced chromosome translocations during meiotic chromosomal segregation. This can lead to infertility, miscarriages or children with abnormalities. Genetic counseling and genetic testing are often offered to families that may carry a translocation. Most balanced translocation carriers are healthy and do not have any symptoms. It is important to distinguish between chromosomal translocations that occur in germ cells, due to errors in meiosis (i.e. during gametogenesis), and those that occur in somatic cells, due to errors in mitosis. The former results in a chromosomal abnormality featured in all cells of the offspring, as in translocation carriers. Somatic translocations, on the other hand, result in abnormalities featured only in the affected cell and its progenitors, as in chronic myelogenous leukemia with the Philadelphia chromosome translocation. Nonreciprocal translocation Nonreciprocal translocation involves the one-way transfer of genes from one chromosome to another nonhomologous chromosome. Robertsonian translocations Robertsonian translocation is a type of translocation caused by breaks at or near the centromeres of two acrocentric chromosomes. The reciprocal exchange of parts gives rise to one large metacentric chromosome and one extremely small chromosome that may be lost from the organism with little effect because it contains few genes. The resulting karyotype in humans leaves only 45 chromosomes, since two chromosomes have fused together. This has no direct effect on the phenotype, since the only genes on the short arms of acrocentrics are common to all of them and are present in variable copy number (nucleolar organiser genes). Robertsonian translocations have been seen involving all combinations of acrocentric chromosomes. The most common translocation in humans involves chromosomes 13 and 14 and is seen in about 0.97 / 1000 newborns. Carriers of Robertsonian translocations are not associated with any phenotypic abnormalities, but there is a risk of unbalanced gametes that lead to miscarriages or abnormal offspring. For example, carriers of Robertsonian translocations involving chromosome 21 have a higher risk of having a child with Down syndrome. This is known as a 'translocation Downs'. This is due to a mis-segregation (nondisjunction) during gametogenesis. The mother has a higher (10%) risk of transmission than the father (1%). Robertsonian translocations involving chromosome 14 also carry a slight risk of uniparental disomy 14 due to trisomy rescue. Role in disease Some human diseases caused by translocations are: Cancer: Several forms of cancer are caused by acquired translocations (as opposed to those present from conception); this has been described mainly in leukemia (acute myelogenous leukemia and chronic myelogenous leukemia). Translocations have also been described in solid malignancies such as Ewing's sarcoma. Infertility: One of the would-be parents carries a balanced translocation, where the parent is asymptomatic but conceived fetuses are not viable. Down syndrome is caused in a minority (5% or less) of cases by a Robertsonian translocation of the chromosome 21 long arm onto the long arm of chromosome 14. Chromosomal translocations between the sex chromosomes can also result in a number of genetic conditions, such as XX male syndrome: caused by a translocation of the SRY gene from the Y to the X chromosome By chromosome Denotation The International System for Human Cytogenetic Nomenclature (ISCN) is used to denote a translocation between chromosomes. The designation t(A;B)(p1;q2) is used to denote a translocation between chromosome A and chromosome B. The information in the second set of parentheses, when given, gives the precise location within the chromosome for chromosomes A and B respectively—with p indicating the short arm of the chromosome, q indicating the long arm, and the numbers after p or q refers to regions, bands and sub-bands seen when staining the chromosome with a staining dye.	Biology and health sciences	Genetics	Biology
561853	https://en.wikipedia.org/wiki/Umber	Umber	Umber is a natural earth pigment consisting of iron oxide and manganese oxide; it has a brownish color that can vary among shades of yellow, red, and green. Umber is considered one of the oldest pigments known to humans, first seen in Ajanta Caves in 200 BC – 600 AD. Umber's advantages are its highly versatile color, warm tone, and quick drying abilities. While some sources indicate that umber's name comes from its geographic origin in Umbria, other scholars suggest that it derives from the Latin word umbra, which means "shadow". The belief that its name derives from the word for shadow is fitting, as the color helps create shadows. The color is primarily produced in Cyprus. Umber is typically mined from open pits or underground mines and ground into a fine powder that is washed to remove impurities. In the 20th century, the rise of synthetic dyes decreased the demand for natural pigments such as umber. History The earliest documented uses of umber date from between 200 BC and 600 AD in the Ajanta Caves found in India. Ocher, a family of earth pigments which includes umber, has been identified in the caves of Altamira in Spain and the Lascaux Cave in France. Some sources indicate that umber was not frequently used in medieval art because of its emphasis on bright and vivid colors. Other sources indicate, however, that umber was used in the Middle Ages to create different shades of brown, most often seen for skin tones. Umber's use in Europe increased in the late 15th century. Umber became more popular during the Renaissance when its versatility, earthy appearance, availability, and inexpensiveness were recognized. Umber gained widespread popularity in Dutch landscape painting in the eighteenth century. Artists recognized the value of umber's high stability, inertness, and drying abilities. It became a standard color within eighteenth-century palettes throughout Europe. Umber's popularity grew during the Baroque period with the rise of the chiaroscuro style. Umber allowed painters to create an intense light and dark contrast. Underpainting was another popular technique for painting that used umber as a base color. Umber was valuable in deploying this technique, creating a range of earth like tones with various layering of color. Toward the end of the 19th century, the Impressionist movement started to use cheaper and more readily available synthetic dyes and reject natural pigments like umber to create mixed hues of brown. The Impressionists chose to make their own browns from mixtures of red, yellow, green, blue and other pigments, particularly the new synthetic pigments such as cobalt blue and emerald green that had just been introduced. In the 20th century, natural umber pigments began to be replaced by pigments made with synthetic iron oxide and manganese oxide. Criticism Beginning in the 17th century, umber was increasingly criticized within the art community. British painter Edward Norgate, prominent with British royalty and aristocracy, called umber "a foul and greasy color." In the 18th century, Spanish painter Antonio Palomino called umber "very false." Jan Blockx, a Belgian painter, opined, "umber should not appear on the palette of the conscientious painter." Visual properties Umber is a natural brown pigment extracted from clay containing iron, manganese, and hydroxides. Umber has diverse hues, ranging from yellow-brown to reddish-brown and even green-brown. The color shade varies depending on the proportions of the components. When heated, umber becomes a more intense color and can look almost black. Burnt umber is produced by calcining the raw version. The raw form of umber is typically used for ceramics because it is less expensive. These warm and earthy tones make it a valuable and versatile pigment for oil painting and other artwork. Umber's high opacity and reactivity of light allow the pigment to have strong hiding power. It is insoluble in water, resistant to alkalis and weak acids, and non-reactive with cement, solvents, oils, and most resins. Umber is known for its stability. Notable occurrences Umber became widely used throughout the Renaissance period for oil paintings. In the Mona Lisa, Leonardo da Vinci used umber for the brown tones throughout his subject’s hair and clothing. Da Vinci also extensively used umber in his painting the Last Supper to create shadows and outlines of the figures. Throughout the Baroque period, many renowned painters used umber. Varieties Raw umber This is the color raw umber. Burnt umber Burnt umber is made by heating raw umber, which dehydrates the iron oxides and changes them partially to the more reddish hematite. It is used for both oil and water color paint. The first recorded use of burnt umber as a color name in English was in 1650.	Physical sciences	Minerals	Earth science
561873	https://en.wikipedia.org/wiki/Earth%20pigment	Earth pigment	Earth pigments are naturally occurring minerals that have been used since prehistoric times as pigments. Among the primary types of earth pigments include the reddish-brown ochres, siennas, and umbers, which contain various amounts of iron oxides and manganese oxides. Other earth pigments include the green earth pigments or , blue earth pigments such as vivianite-based "blue ochre", white earth pigments such as chalk, and black earth pigments such as charcoal. Earth pigments are known for their fast drying time in oil painting, relative inexpensiveness, and lightfastness. Cave paintings done in sienna still survive today. Production After mining, the mineral used for making a pigment is ground to a very fine powder (if not already in the form of clay), washed to remove water-soluble components, dried, and ground again to powder. For some pigments, notably sienna and umber, the color can be deepened by heating (calcination) in a process known as "burning", although it does not involve oxidation but instead dehydration.	Physical sciences	Minerals	Earth science
562007	https://en.wikipedia.org/wiki/Magnesium%20chloride	Magnesium chloride	Magnesium chloride is an inorganic compound with the formula . It forms hydrates , where n can range from 1 to 12. These salts are colorless or white solids that are highly soluble in water. These compounds and their solutions, both of which occur in nature, have a variety of practical uses. Anhydrous magnesium chloride is the principal precursor to magnesium metal, which is produced on a large scale. Hydrated magnesium chloride is the form most readily available. Production Magnesium chloride can be extracted from brine or sea water. In North America and South America too, it is produced primarily from Great Salt Lake brine. In the Jordan Valley, it is obtained from the Dead Sea. The mineral bischofite () is extracted (by solution mining) out of ancient seabeds, for example, the Zechstein seabed in northwest Europe. Some deposits result from high content of magnesium chloride in the primordial ocean. Some magnesium chloride is made from evaporation of seawater. In the Dow process, magnesium chloride is regenerated from magnesium hydroxide using hydrochloric acid: It can also be prepared from magnesium carbonate by a similar reaction. Structure crystallizes in the cadmium chloride motif, therefore it loses water upon heating: n = 12 (−16.4 °C), 8 (−3.4 °C), 6 (116.7 °C), 4 (181 °C), 2 (about 300 °C). In the hexahydrate, the is also octahedral, being coordinated to six water ligands. The octahydrate and the dodecahydrate can be crystallized from water below 298K. As verified by X-ray crystallography, these "higher" hydrates also feature [Mg(H2O)6]2+ ions. A decahydrate has also been crystallized. Preparation, general properties Anhydrous is produced industrially by heating the complex salt named hexamminemagnesium dichloride . The thermal dehydration of the hydrates (n = 6, 12) does not occur straightforwardly. As suggested by the existence of hydrates, anhydrous is a Lewis acid, although a weak one. One derivative is tetraethylammonium tetrachloromagnesate . The adduct is another. In the coordination polymer with the formula , Mg adopts an octahedral geometry. The Lewis acidity of magnesium chloride is reflected in its deliquescence, meaning that it attracts moisture from the air to the extent that the solid turns into a liquid. Applications Precursor to metallic magnesium Anhydrous is the main precursor to metallic magnesium. The reduction of into metallic Mg is performed by electrolysis in molten salt. As it is also the case for aluminium, an electrolysis in aqueous solution is not possible as the produced metallic magnesium would immediately react with water, or in other words that the water would be reduced into gaseous before Mg reduction could occur. So, the direct electrolysis of molten in the absence of water is required because the reduction potential to obtain Mg is lower than the stability domain of water on an Eh–pH diagram (Pourbaix diagram). The production of metallic magnesium at the cathode (reduction reaction) is accompanied by the oxidation of the chloride anions at the anode with release of gaseous chlorine. This process is developed at a large industrial scale. Dust and erosion control Magnesium chloride is one of many substances used for dust control, soil stabilization, and wind erosion mitigation. When magnesium chloride is applied to roads and bare soil areas, both positive and negative performance issues occur which are related to many application factors. Catalysis Ziegler-Natta catalysts, used commercially to produce polyolefins, often contain as a catalyst support. The introduction of supports increases the activity of traditional catalysts and allowed the development of highly stereospecific catalysts for the production of polypropylene. Magnesium chloride is also a Lewis acid catalyst in aldol reactions. Ice control Magnesium chloride is used for low-temperature de-icing of highways, sidewalks, and parking lots. When highways are treacherous due to icy conditions, magnesium chloride is applied to help prevent ice from bonding to the pavement, allowing snow plows to clear treated roads more efficiently. For the purpose of preventing ice from forming on pavement, magnesium chloride is applied in three ways: anti-icing, which involves spreading it on roads to prevent snow from sticking and forming; prewetting, which means a liquid formulation of magnesium chloride is sprayed directly onto salt as it is being spread onto roadway pavement, wetting the salt so that it sticks to the road; and pretreating, when magnesium chloride and salt are mixed together before they are loaded onto trucks and spread onto paved roads. Calcium chloride damages concrete twice as fast as magnesium chloride. The amount of magnesium chloride is supposed to be controlled when it is used for de-icing as it may cause pollution to the environment. Nutrition and medicine Magnesium chloride is used in nutraceutical and pharmaceutical preparations. The hexahydrate is sometimes advertised as "magnesium oil". Cuisine Magnesium chloride (E511) is an important coagulant used in the preparation of tofu from soy milk. In Japan it is sold as nigari (にがり, derived from the Japanese word for "bitter"), a white powder produced from seawater after the sodium chloride has been removed, and the water evaporated. In China, it is called lushui (卤水). Nigari or Iushui is, in fact, natural magnesium chloride, meaning that it is not completely refined (it contains up to 5% magnesium sulfate and various minerals). The crystals originate from lakes in the Chinese province of Qinghai, to be then reworked in Japan. Gardening and horticulture Because magnesium is a mobile nutrient, magnesium chloride can be effectively used as a substitute for magnesium sulfate (Epsom salt) to help correct magnesium deficiency in plants via foliar feeding. The recommended dose of magnesium chloride is smaller than the recommended dose of magnesium sulfate (20 g/L). This is due primarily to the chlorine present in magnesium chloride, which can easily reach toxic levels if over-applied or applied too often. It has been found that higher concentrations of magnesium in tomato and some pepper plants can make them more susceptible to disease caused by infection of the bacterium Xanthomonas campestris, since magnesium is essential for bacterial growth. Wastewater treatment It is used to supply the magnesium necessary to precipitate phosphorus in the form of struvite from agricultural waste as well as human urine. Occurrence Magnesium concentrations in natural seawater are between 1250 and 1350 mg/L, around 3.7% of the total seawater mineral content. Dead Sea minerals contain a significantly higher magnesium chloride ratio, 50.8%. Carbonates and calcium are essential for all growth of corals, coralline algae, clams, and invertebrates. Magnesium can be depleted by mangrove plants and the use of excessive limewater or by going beyond natural calcium, alkalinity, and pH values. The most common mineral form of magnesium chloride is its hexahydrate, bischofite. Anhydrous compound occurs very rarely, as chloromagnesite. Magnesium chloride-hydroxides, korshunovskite and nepskoeite, are also very rare. Toxicology Magnesium ions are bitter-tasting, and magnesium chloride solutions are bitter in varying degrees, depending on the concentration. Magnesium toxicity from magnesium salts is rare in healthy individuals with a normal diet, because excess magnesium is readily excreted in urine by the kidneys. A few cases of oral magnesium toxicity have been described in persons with normal renal function ingesting large amounts of magnesium salts, but it is rare. If a large amount of magnesium chloride is eaten, it will have effects similar to magnesium sulfate, causing diarrhea, although the sulfate also contributes to the laxative effect in magnesium sulfate, so the effect from the chloride is not as severe. Plant toxicity Chloride () and magnesium () are both essential nutrients important for normal plant growth. Too much of either nutrient may harm a plant, although foliar chloride concentrations are more strongly related with foliar damage than magnesium. High concentrations of ions in the soil may be toxic or change water relationships such that the plant cannot easily accumulate water and nutrients. Once inside the plant, chloride moves through the water-conducting system and accumulates at the margins of leaves or needles, where dieback occurs first. Leaves are weakened or killed, which can lead to the death of the tree.	Physical sciences	Halide salts	Chemistry
562067	https://en.wikipedia.org/wiki/Brillouin%20zone	Brillouin zone	In mathematics and solid state physics, the first Brillouin zone (named after Léon Brillouin) is a uniquely defined primitive cell in reciprocal space. In the same way the Bravais lattice is divided up into Wigner–Seitz cells in the real lattice, the reciprocal lattice is broken up into Brillouin zones. The boundaries of this cell are given by planes related to points on the reciprocal lattice. The importance of the Brillouin zone stems from the description of waves in a periodic medium given by Bloch's theorem, in which it is found that the solutions can be completely characterized by their behavior in a single Brillouin zone. The first Brillouin zone is the locus of points in reciprocal space that are closer to the origin of the reciprocal lattice than they are to any other reciprocal lattice points (see the derivation of the Wigner–Seitz cell). Another definition is as the set of points in k-space that can be reached from the origin without crossing any Bragg plane. Equivalently, this is the Voronoi cell around the origin of the reciprocal lattice. There are also second, third, etc., Brillouin zones, corresponding to a sequence of disjoint regions (all with the same volume) at increasing distances from the origin, but these are used less frequently. As a result, the first Brillouin zone is often called simply the Brillouin zone. In general, the n-th Brillouin zone consists of the set of points that can be reached from the origin by crossing exactly n − 1 distinct Bragg planes. A related concept is that of the irreducible Brillouin zone, which is the first Brillouin zone reduced by all of the symmetries in the point group of the lattice (point group of the crystal). The concept of a Brillouin zone was developed by Léon Brillouin (1889–1969), a French physicist. Within the Brillouin zone, a constant-energy surface represents the loci of all the -points (that is, all the electron momentum values) that have the same energy. Fermi surface is a special constant-energy surface that separates the unfilled orbitals from the filled ones at zero kelvin. Critical points Several points of high symmetry are of special interest – these are called critical points. Other lattices have different types of high-symmetry points. They can be found in the illustrations below.	Physical sciences	Crystallography	Physics
7070301	https://en.wikipedia.org/wiki/Telescope	Telescope	A telescope is a device used to observe distant objects by their emission, absorption, or reflection of electromagnetic radiation. Originally, it was an optical instrument using lenses, curved mirrors, or a combination of both to observe distant objects – an optical telescope. Nowadays, the word "telescope" is defined as a wide range of instruments capable of detecting different regions of the electromagnetic spectrum, and in some cases other types of detectors. The first known practical telescopes were refracting telescopes with glass lenses and were invented in the Netherlands at the beginning of the 17th century. They were used for both terrestrial applications and astronomy. The reflecting telescope, which uses mirrors to collect and focus light, was invented within a few decades of the first refracting telescope. In the 20th century, many new types of telescopes were invented, including radio telescopes in the 1930s and infrared telescopes in the 1960s. Etymology The word telescope was coined in 1611 by the Greek mathematician Giovanni Demisiani for one of Galileo Galilei's instruments presented at a banquet at the Accademia dei Lincei. In the Starry Messenger, Galileo had used the Latin term . The root of the word is from the Ancient Greek τῆλε, tele 'far' and σκοπεῖν, skopein 'to look or see'; τηλεσκόπος, teleskopos 'far-seeing'. History The earliest existing record of a telescope was a 1608 patent submitted to the government in the Netherlands by Middelburg spectacle maker Hans Lipperhey for a refracting telescope. The actual inventor is unknown but word of it spread through Europe. Galileo heard about it and, in 1609, built his own version, and made his telescopic observations of celestial objects. The idea that the objective, or light-gathering element, could be a mirror instead of a lens was being investigated soon after the invention of the refracting telescope. The potential advantages of using parabolic mirrors—reduction of spherical aberration and no chromatic aberration—led to many proposed designs and several attempts to build reflecting telescopes. In 1668, Isaac Newton built the first practical reflecting telescope, of a design which now bears his name, the Newtonian reflector. The invention of the achromatic lens in 1733 partially corrected color aberrations present in the simple lens and enabled the construction of shorter, more functional refracting telescopes. Reflecting telescopes, though not limited by the color problems seen in refractors, were hampered by the use of fast tarnishing speculum metal mirrors employed during the 18th and early 19th century—a problem alleviated by the introduction of silver coated glass mirrors in 1857, and aluminized mirrors in 1932. The maximum physical size limit for refracting telescopes is about , dictating that the vast majority of large optical researching telescopes built since the turn of the 20th century have been reflectors. The largest reflecting telescopes currently have objectives larger than , and work is underway on several 30–40m designs. The 20th century also saw the development of telescopes that worked in a wide range of wavelengths from radio to gamma-rays. The first purpose-built radio telescope went into operation in 1937. Since then, a large variety of complex astronomical instruments have been developed. In space Since the atmosphere is opaque for most of the electromagnetic spectrum, only a few bands can be observed from the Earth's surface. These bands are visible – near-infrared and a portion of the radio-wave part of the spectrum. For this reason there are no X-ray or far-infrared ground-based telescopes as these have to be observed from orbit. Even if a wavelength is observable from the ground, it might still be advantageous to place a telescope on a satellite due to issues such as clouds, astronomical seeing and light pollution. The disadvantages of launching a space telescope include cost, size, maintainability and upgradability. Some examples of space telescopes from NASA are the Hubble Space Telescope that detects visible light, ultraviolet, and near-infrared wavelengths, the Spitzer Space Telescope that detects infrared radiation, and the Kepler Space Telescope that discovered thousands of exoplanets. The latest telescope that was launched was the James Webb Space Telescope on December 25, 2021, in Kourou, French Guiana. The Webb telescope detects infrared light. By electromagnetic spectrum The name "telescope" covers a wide range of instruments. Most detect electromagnetic radiation, but there are major differences in how astronomers must go about collecting light (electromagnetic radiation) in different frequency bands. As wavelengths become longer, it becomes easier to use antenna technology to interact with electromagnetic radiation (although it is possible to make very tiny antenna). The near-infrared can be collected much like visible light; however, in the far-infrared and submillimetre range, telescopes can operate more like a radio telescope. For example, the James Clerk Maxwell Telescope observes from wavelengths from 3 μm (0.003 mm) to 2000 μm (2 mm), but uses a parabolic aluminum antenna. On the other hand, the Spitzer Space Telescope, observing from about 3 μm (0.003 mm) to 180 μm (0.18 mm) uses a mirror (reflecting optics). Also using reflecting optics, the Hubble Space Telescope with Wide Field Camera 3 can observe in the frequency range from about 0.2 μm (0.0002 mm) to 1.7 μm (0.0017 mm) (from ultra-violet to infrared light). With photons of the shorter wavelengths, with the higher frequencies, glancing-incident optics, rather than fully reflecting optics are used. Telescopes such as TRACE and SOHO use special mirrors to reflect extreme ultraviolet, producing higher resolution and brighter images than are otherwise possible. A larger aperture does not just mean that more light is collected, it also enables a finer angular resolution. Telescopes may also be classified by location: ground telescope, space telescope, or flying telescope. They may also be classified by whether they are operated by professional astronomers or amateur astronomers. A vehicle or permanent campus containing one or more telescopes or other instruments is called an observatory. Radio and submillimeter Radio telescopes are directional radio antennas that typically employ a large dish to collect radio waves. The dishes are sometimes constructed of a conductive wire mesh whose openings are smaller than the wavelength being observed. Unlike an optical telescope, which produces a magnified image of the patch of sky being observed, a traditional radio telescope dish contains a single receiver and records a single time-varying signal characteristic of the observed region; this signal may be sampled at various frequencies. In some newer radio telescope designs, a single dish contains an array of several receivers; this is known as a focal-plane array. By collecting and correlating signals simultaneously received by several dishes, high-resolution images can be computed. Such multi-dish arrays are known as astronomical interferometers and the technique is called aperture synthesis. The 'virtual' apertures of these arrays are similar in size to the distance between the telescopes. As of 2005, the record array size is many times the diameter of the Earth – using space-based very-long-baseline interferometry (VLBI) telescopes such as the Japanese HALCA (Highly Advanced Laboratory for Communications and Astronomy) VSOP (VLBI Space Observatory Program) satellite. Aperture synthesis is now also being applied to optical telescopes using optical interferometers (arrays of optical telescopes) and aperture masking interferometry at single reflecting telescopes. Radio telescopes are also used to collect microwave radiation, which has the advantage of being able to pass through the atmosphere and interstellar gas and dust clouds. Some radio telescopes such as the Allen Telescope Array are used by programs such as SETI and the Arecibo Observatory to search for extraterrestrial life. Infrared Visible light An optical telescope gathers and focuses light mainly from the visible part of the electromagnetic spectrum. Optical telescopes increase the apparent angular size of distant objects as well as their apparent brightness. For the image to be observed, photographed, studied, and sent to a computer, telescopes work by employing one or more curved optical elements, usually made from glass lenses and/or mirrors, to gather light and other electromagnetic radiation to bring that light or radiation to a focal point. Optical telescopes are used for astronomy and in many non-astronomical instruments, including: theodolites (including transits), spotting scopes, monoculars, binoculars, camera lenses, and spyglasses. There are three main optical types: The refracting telescope which uses lenses to form an image. The reflecting telescope which uses an arrangement of mirrors to form an image. The catadioptric telescope which uses mirrors combined with lenses to form an image. A Fresnel imager is a proposed ultra-lightweight design for a space telescope that uses a Fresnel lens to focus light. Beyond these basic optical types there are many sub-types of varying optical design classified by the task they perform such as astrographs, comet seekers and solar telescopes. Ultraviolet Most ultraviolet light is absorbed by the Earth's atmosphere, so observations at these wavelengths must be performed from the upper atmosphere or from space. X-ray X-rays are much harder to collect and focus than electromagnetic radiation of longer wavelengths. X-ray telescopes can use X-ray optics, such as Wolter telescopes composed of ring-shaped 'glancing' mirrors made of heavy metals that are able to reflect the rays just a few degrees. The mirrors are usually a section of a rotated parabola and a hyperbola, or ellipse. In 1952, Hans Wolter outlined 3 ways a telescope could be built using only this kind of mirror. Examples of space observatories using this type of telescope are the Einstein Observatory, ROSAT, and the Chandra X-ray Observatory. In 2012 the NuSTAR X-ray Telescope was launched which uses Wolter telescope design optics at the end of a long deployable mast to enable photon energies of 79 keV. Gamma ray Higher energy X-ray and gamma ray telescopes refrain from focusing completely and use coded aperture masks: the patterns of the shadow the mask creates can be reconstructed to form an image. X-ray and Gamma-ray telescopes are usually installed on high-flying balloons or Earth-orbiting satellites since the Earth's atmosphere is opaque to this part of the electromagnetic spectrum. An example of this type of telescope is the Fermi Gamma-ray Space Telescope which was launched in June 2008. The detection of very high energy gamma rays, with shorter wavelength and higher frequency than regular gamma rays, requires further specialization. Such detections can be made either with the Imaging Atmospheric Cherenkov Telescopes (IACTs) or with Water Cherenkov Detectors (WCDs). Examples of IACTs are H.E.S.S. and VERITAS with the next-generation gamma-ray telescope, the Cherenkov Telescope Array (CTA), currently under construction. HAWC and LHAASO are examples of gamma-ray detectors based on the Water Cherenkov Detectors. A discovery in 2012 may allow focusing gamma-ray telescopes. At photon energies greater than 700 keV, the index of refraction starts to increase again. Lists of telescopes List of optical telescopes List of largest optical reflecting telescopes List of largest optical refracting telescopes List of largest optical telescopes historically List of radio telescopes List of solar telescopes List of space observatories List of telescope parts and construction List of telescope types	Technology	Optical	null
7071096	https://en.wikipedia.org/wiki/Engineering%20design%20process	Engineering design process	The engineering design process, also known as the engineering method, is a common series of steps that engineers use in creating functional products and processes. The process is highly iterative – parts of the process often need to be repeated many times before another can be entered – though the part(s) that get iterated and the number of such cycles in any given project may vary. It is a decision making process (often iterative) in which the engineering sciences, basic sciences and mathematics are applied to convert resources optimally to meet a stated objective. Among the fundamental elements of the design process are the establishment of objectives and criteria, synthesis, analysis, construction, testing and evaluation. Common stages of the engineering design process It's important to understand that there are various framings/articulations of the engineering design process. Different terminology employed may have varying degrees of overlap, which affects what steps get stated explicitly or deemed "high level" versus subordinate in any given model. This, of course, applies as much to any particular example steps/sequences given here. One example framing of the engineering design process delineates the following stages: research, conceptualization, feasibility assessment, establishing design requirements, preliminary design, detailed design, production planning and tool design, and production. Others, noting that "different authors (in both research literature and in textbooks) define different phases of the design process with varying activities occurring within them," have suggested more simplified/generalized models – such as problem definition, conceptual design, preliminary design, detailed design, and design communication. Another summary of the process, from European engineering design literature, includes clarification of the task, conceptual design, embodiment design, detail design. (NOTE: In these examples, other key aspects – such as concept evaluation and prototyping – are subsets and/or extensions of one or more of the listed steps.) Research Various stages of the design process (and even earlier) can involve a significant amount of time spent on locating information and research. Consideration should be given to the existing applicable literature, problems and successes associated with existing solutions, costs, and marketplace needs. The source of information should be relevant. Reverse engineering can be an effective technique if other solutions are available on the market. Other sources of information include the Internet, local libraries, available government documents, personal organizations, trade journals, vendor catalogs and individual experts available. Design requirements Establishing design requirements and conducting requirement analysis, sometimes termed problem definition (or deemed a related activity), is one of the most important elements in the design process in certain industries, and this task is often performed at the same time as a feasibility analysis. The design requirements control the design of the product or process being developed, throughout the engineering design process. These include basic things like the functions, attributes, and specifications – determined after assessing user needs. Some design requirements include hardware and software parameters, maintainability, availability, and testability. Feasibility In some cases, a feasibility study is carried out after which schedules, resource plans and estimates for the next phase are developed. The feasibility study is an evaluation and analysis of the potential of a proposed project to support the process of decision making. It outlines and analyses alternatives or methods of achieving the desired outcome. The feasibility study helps to narrow the scope of the project to identify the best scenario. A feasibility report is generated following which Post Feasibility Review is performed. The purpose of a feasibility assessment is to determine whether the engineer's project can proceed into the design phase. This is based on two criteria: the project needs to be based on an achievable idea, and it needs to be within cost constraints. It is important to have engineers with experience and good judgment to be involved in this portion of the feasibility study. Concept generation A concept study (conceptualization, conceptual design) is often a phase of project planning that includes producing ideas and taking into account the pros and cons of implementing those ideas. This stage of a project is done to minimize the likelihood of error, manage costs, assess risks, and evaluate the potential success of the intended project. In any event, once an engineering issue or problem is defined, potential solutions must be identified. These solutions can be found by using ideation, the mental process by which ideas are generated. In fact, this step is often termed Ideation or "Concept Generation." The following are widely used techniques: trigger word – a word or phrase associated with the issue at hand is stated, and subsequent words and phrases are evoked. morphological analysis – independent design characteristics are listed in a chart, and different engineering solutions are proposed for each solution. Normally, a preliminary sketch and short report accompany the morphological chart. synectics – the engineer imagines him or herself as the item and asks, "What would I do if I were the system?" This unconventional method of thinking may find a solution to the problem at hand. The vital aspects of the conceptualization step is synthesis. Synthesis is the process of taking the element of the concept and arranging them in the proper way. Synthesis creative process is present in every design. brainstorming – this popular method involves thinking of different ideas, typically as part of a small group, and adopting these ideas in some form as a solution to the problem Various generated ideas must then undergo a concept evaluation step, which utilizes various tools to compare and contrast the relative strengths and weakness of possible alternatives. Preliminary design The preliminary design, or high-level design includes (also called FEED or Basic design), often bridges a gap between design conception and detailed design, particularly in cases where the level of conceptualization achieved during ideation is not sufficient for full evaluation. So in this task, the overall system configuration is defined, and schematics, diagrams, and layouts of the project may provide early project configuration. (This notably varies a lot by field, industry, and product.) During detailed design and optimization, the parameters of the part being created will change, but the preliminary design focuses on creating the general framework to build the project on. S. Blanchard and J. Fabrycky describe it as: “The ‘whats’ initiating conceptual design produce ‘hows’ from the conceptual design evaluation effort applied to feasible conceptual design concepts. Next, the ‘hows’ are taken into preliminary design through the means of allocated requirements. There they become ‘whats’ and drive preliminary design to address ‘hows’ at this lower level.” Detailed design Following FEED is the Detailed Design (Detailed Engineering) phase, which may consist of procurement of materials as well. This phase further elaborates each aspect of the project/product by complete description through solid modeling, drawings as well as specifications. Computer-aided design (CAD) programs have made the detailed design phase more efficient. For example, a CAD program can provide optimization to reduce volume without hindering a part's quality. It can also calculate stress and displacement using the finite element method to determine stresses throughout the part. Production planning The production planning and tool design consists of planning how to mass-produce the product and which tools should be used in the manufacturing process. Tasks to complete in this step include selecting materials, selection of the production processes, determination of the sequence of operations, and selection of tools such as jigs, fixtures, metal cutting and metal or plastics forming tools. This task also involves additional prototype testing iterations to ensure the mass-produced version meets qualification testing standards. Comparison with the scientific method Engineering is formulating a problem that can be solved through design. Science is formulating a question that can be solved through investigation. The engineering design process bears some similarity to the scientific method. Both processes begin with existing knowledge, and gradually become more specific in the search for knowledge (in the case of "pure" or basic science) or a solution (in the case of "applied" science, such as engineering). The key difference between the engineering process and the scientific process is that the engineering process focuses on design, creativity and innovation while the scientific process emphasizes explanation, prediction and discovery (observation). Degree programs Methods are being taught and developed in Universities including: Engineering Design, University of Bristol Faculty of Engineering Dyson School of Design Engineering, Imperial College London TU Delft, Industrial Design Engineering. University of Waterloo, Systems Design Engineering	Technology	Basics	null
6914315	https://en.wikipedia.org/wiki/Maja%20squinado	Maja squinado	Maja squinado (the European spider crab, spiny spider crab or spinous spider crab) is a species of migratory crab found in the Mediterranean Sea. The appearance of the European spider crab is similar to the much larger Japanese spider crab, although the European spider crab belongs to the family Majidae, and the Japanese spider crab belongs to a different family of crabs, the Macrocheiridae. Young The young of M. squinado are slightly longer than 1mm after hatching, and weigh approximately 0.12 mg at this time. Within 4–8 days, the larva moults numerous times, finally ending with morphological changes that presumably include the further development and increase in size of the cephalothorax. In a second phase, the Carapace grows to a length of approx. 2mm, and weighs approx. 0.3 mg. The larva then undergoes metamorphosis to the first juvenile instar, and changes its planktonic life to a benthic one (living on the sea floor). Its appearance is also similar to that of the adult animal. From this stage only growth and the formation of sexual maturity follows. In observations under laboratory conditions, approx. 10.5% of the hatched zoea made it to this stage. The same conditions in terms of food, temperature and the like cannot be created in a laboratory. Animals in the first juvenile stage perform their first moult about 21 days after hatching, and therefore enter the second juvenile stage. Here there is a considerable increase in the length of the carapace to approx. 4.51 mm. The second moult marks the beginning of the third juvenile stage, the animal now has the appearance of the adult, with a carapace length of approx. 5.63 mm, but is not sexually mature. Juvenile animals spend another 2 years moulting and growing in size. The juvenile animals live in shallow water in winter, between rocks in coastal kelp forests. They spend the summer on small rocky reefs at a depth of only about 4 m. After this time, they reach a carapace length between 6–13 cm, with no noticeable sex-specific differences. During this time they are not yet sexually mature. There are two main periods for the critical moults that follow the approximately two-year period of growth leading to sexual maturity: the first, the prepubertal, in April, and the second, the pubertal, from July to October. However, in captive animals it has been noticed that in very large individuals that are in the phase before one of the two moults, one moult may be lost entirely, or be very late. Likewise, three moults have been observed on some individual specimens. The average time interval between the two critical moults is 104 days. Typically, the carapace length in animals that are already comparatively large increases less after moulting, relative to the initial size, than that of smaller animals. This also explains why there is a smaller increase in length (approx. 27%) in the pubescent moult than in the prepubertal (approx. 36%). Behavior Migrations generally take place in autumn, with some crabs covering over in 8 months. All crabs are vulnerable to predation when moulting, and M. squinado becomes gregarious around that time, presumably for defense against predators. Females can produce up to four broods per year. M. squinado has been documented to feed on macroalgae and benthic invertebrates. From a 1992 study done in Galicia, seaweeds from the Laminariaceae, Corallina spp., molluscs, the gastropods Bittium spp., Trochiidae, the bivalve Mytilus spp., echinoderms, and others were observed as part of the diet of this particular species. Fishery M. squinado is the subject of commercial fishery, with over 5,000 tonnes caught annually, more than 70% the coast of France, over 10% off the coast of the United Kingdom, 6% from the Channel Islands, 3% from each of Spain and Ireland, 2% from Croatia, 1% from Portugal, and the remainder coming from Montenegro, Denmark, and Morocco, although official production figures are open to doubt. The European Union imposes a minimum landing size of 120 mm for M. squinado, and some individual countries have other regulations, such as a ban on landing egg-bearing females in Spain and a closed season in France and the Channel Islands. Taxonomy A review of the species complex around M. squinado was able to differentiate between specimens from the Mediterranean Sea and those from the Atlantic, and concluded that the Atlantic specimens were a separate species, called Maja brachydactyla Balss, 1922. The specific epithet squinado derives from the Provençal name for the species – , , or — recorded by Rondelet as early as 1554.	Biology and health sciences	Crabs and hermit crabs	Animals
6915193	https://en.wikipedia.org/wiki/Stereospondyli	Stereospondyli	The Stereospondyli are a group of extinct temnospondyl amphibians that existed primarily during the Mesozoic period. They are known from all seven continents and were common components of many Triassic ecosystems, likely filling a similar ecological niche to modern crocodilians prior to the diversification of pseudosuchian archosaurs. Classification and anatomy The group was first defined by Zittel (1888) on the recognition of the distinctive vertebral anatomy of the best known stereospondyls of the time, such as Mastodonsaurus and Metoposaurus. The term 'stereospondylous' as a descriptor of vertebral anatomy was coined the following year by Fraas, referring to a vertebral position consisting largely or entirely of the intercentrum in addition to the neural arch. While the name 'Stereospondyli' is derived from the stereospondylous vertebral condition, there is a diversity of vertebral morphologies among stereospondyls, including the diplospondylous ('tupilakosaurid') condition, where the arch sits between the corresponding intercentrum and pleurocentrum, and the plagiosaurid condition, where a single large centrum ossification (identity unknown) is present, and the arch sits between subsequent vertebral positions. The concept of Stereospondyli has thus undergone repeated and frequent revisions by different workers. Defining features include a tight articulation between the parasphenoid and the pterygoid and a stapedial groove. Evolutionary history Stereospondyls first definitively appeared during the early Permian, as represented by fragmentary remains of a rhinesuchid from the Pedra de Fogo Formation of Brazil. Rhinesuchids are one of the earliest groups of stereospondyls to appear in the fossil record and are predominantly a late Permian clade, with only one species, Broomistega putterilli, from the Early Triassic of South Africa. However, almost all other groups of stereospondyls are not known from any Paleozoic deposits, which remained dominated by non-stereospondyl stereospondylomorphs. The taxonomically unresolved Peltobatrachus pustulatus, which has historically been regarded as a stereospondyl, is also known from the late Permian of Tanzania. Several more fragmentary records are known from horizons spanning the Permo-Triassic boundary in South America, such as the rhinesuchid-like Arachana nigra from Uruguay and an indeterminate mastodonsaurid from Uruguay. Following the Permo-Triassic mass extinction, stereospondyls are abundantly represented in the fossil record, particularly from Russia, South Africa, and Australia. This led Yates & Warren (2000) to propose that stereospondyls had sheltered in a high-latitude refugium that would have been somewhat shielded from the global effects of the extinction, and that they subsequently radiated from present-day Australia or Antarctica. Recent discoveries of a diverse rhinesuchid community in South America alongside non-stereospondyl stereospondylomorphs have led to an alternative hypothesis for a radiation from western Gondwana in South America. By the end of the Early Triassic, virtually all major clades of stereospondyls had appeared in the fossil record, although some were more geographically localized (e.g., lapillopsids, rhytidosteids) than those with cosmopolitan distributions (e.g., capitosauroids, trematosauroids). Stereospondyls were the latest-surviving temnospondyl group. With the diversification of crocodile-like archosaurs and an extinction event at the end of the Triassic, most other temnospondyls disappeared. Chigutisaurid brachyopoids persisted into the Jurassic in Asia and Australia, including Koolasuchus, the youngest known stereospondyl (late Early Cretaceous) from what is now Australia. There is also sparse evidence for the persistence of some trematosauroids into the Jurassic of Asia. If the recent hypothesis that Chinlestegophis, a Late Triassic stereospondyl from North America, is indeed a stem caecilian is correct, then stereospondyls would survive to the present day. Lifestyle and ecology Stereospondyls were particularly diverse during the Early Triassic, with small-bodied taxa such as lapillopsids and lydekkerinids that were likely more terrestrially capable present alongside larger taxa that would continue into the Middle Triassic, such as brachyopoids and trematosauroids. The vast majority of stereospondyls, particularly the large-bodied taxa, have been inferred to have been obligately aquatic based on features of the external anatomy such as a well-developed lateral line system, poorly ossified postcranial skeleton, and occasional preservation of proxies of external gills. Many taxa also reflect adaptations for an aquatic lifestyle as evidence in bone histology, which is pachyostotic in many taxa, although some studies suggest a greater terrestrial ability than historically inferred. Most of the aquatic taxa resided in freshwater environments, but some trematosauroids in particular are thought to have been euryhaline based on their preservation in marine sediments with marine organisms. While stereospondyls are often compared to modern crocodilians, the presence of multiple temnospondyls in some environments and the range of morphologies across Stereospondyli indicates that at least some clades occupied drastically different ecological niches, such as benthic ambush predators. Some groups, such as metoposaurids, are often recovered from large monotaxic bone beds interpreted as evidence of aggregation prior to mass death. Relationships Phylogeny Gallery	Biology and health sciences	Prehistoric amphibians	Animals
11671698	https://en.wikipedia.org/wiki/Soil%20salinity%20control	Soil salinity control	Soil salinity control refers to controlling the process and progress of soil salinity to prevent soil degradation by salination and reclamation of already salty (saline) soils. Soil reclamation is also known as soil improvement, rehabilitation, remediation, recuperation, or amelioration. The primary man-made cause of salinization is irrigation. River water or groundwater used in irrigation contains salts, which remain in the soil after the water has evaporated. The primary method of controlling soil salinity is to permit 10–20% of the irrigation water to leach the soil, so that it will be drained and discharged through an appropriate drainage system. The salt concentration of the drainage water is normally 5 to 10 times higher than that of the irrigation water which meant that salt export will more closely match salt import and it will not accumulate. Problems with soil salinity Salty (saline) soils have high salt content. The predominant salt is normally sodium chloride (NaCl, "table salt"). Saline soils are therefore also sodic soils but there may be sodic soils that are not saline, but alkaline. According to a study by UN University, about , representing 20% of the world's irrigated lands are affected, up from in the early 1990s. In the Indo-Gangetic Plain, home to over 10% of the world's population, crop yield losses for wheat, rice, sugarcane and cotton grown on salt-affected lands could be 40%, 45%, 48%, and 63%, respectively. Salty soils are a common feature and an environmental problem in irrigated lands in arid and semi-arid regions, resulting in poor or little crop production. The causes of salty soils are often associated with high water tables, which are caused by a lack of natural subsurface drainage to the underground. Poor subsurface drainage may be caused by insufficient transport capacity of the aquifer or because water cannot exit the aquifer, for instance, if the aquifer is situated in a topographical depression. Worldwide, the major factor in the development of saline soils is a lack of precipitation. Most naturally saline soils are found in (semi) arid regions and climates of the earth. Primary cause Man-made salinization is primarily caused by salt found in irrigation water. All irrigation water derived from rivers or groundwater, regardless of water purity, contains salts that remain behind in the soil after the water has evaporated. For example, assuming irrigation water with a low salt concentration of 0.3 g/L (equal to 0.3 kg/m3 corresponding to an electric conductivity of about 0.5 FdS/m) and a modest annual supply of irrigation water of 10,000 m3/ha (almost 3 mm/day) brings 3,000 kg salt/ha each year. With the absence of sufficient natural drainage (as in waterlogged soils), and proper leaching and drainage program to remove salts, this would lead to high soil salinity and reduced crop yields in the long run. Much of the water used in irrigation has a higher salt content than 0.3 g/L, compounded by irrigation projects using a far greater annual supply of water. Sugar cane, for example, needs about 20,000 m3/ha of water per year. As a result, irrigated areas often receive more than 3,000 kg/ha of salt per year, with some receiving as much as 10,000 kg/ha/year. Secondary cause The secondary cause of salinization is waterlogging in irrigated land. Irrigation causes changes to the natural water balance of irrigated lands. Large quantities of water in irrigation projects are not consumed by plants and must go somewhere. In irrigation projects, it is impossible to achieve 100% irrigation efficiency where all the irrigation water is consumed by the plants. The maximum attainable irrigation efficiency is about 70%, but usually, it is less than 60%. This means that minimum 30%, but usually more than 40% of the irrigation water is not evaporated and it must go somewhere. Most of the water lost this way is stored underground which can change the original hydrology of local aquifers considerably. Many aquifers cannot absorb and transport these quantities of water, and so the water table rises leading to waterlogging. Waterlogging causes three problems: The shallow water table and lack of oxygenation of the root zone reduces the yield of most crops. It leads to an accumulation of salts brought in with the irrigation water as their removal through the aquifer is blocked. With the upward seepage of groundwater, more salts are brought into the soil and the salination is aggravated. Aquifer conditions in irrigated land and the groundwater flow have an important role in soil salinization, as illustrated here: Salt affected area Normally, the salinization of agricultural land affects a considerable area of 20% to 30% in irrigation projects. When the agriculture in such a fraction of the land is abandoned, a new salt and water balance is attained, a new equilibrium is reached and the situation becomes stable. In India alone, thousands of square kilometers have been severely salinized. China and Pakistan do not lag far behind (perhaps China has even more salt affected land than India). A regional distribution of the 3,230,000 km2 of saline land worldwide is shown in the following table derived from the FAO/UNESCO Soil Map of the World. Spatial variation Although the principles of the processes of salinization are fairly easy to understand, it is more difficult to explain why certain parts of the land suffer from the problems and other parts do not, or to predict accurately which part of the land will fall victim. The main reason for this is the variation of natural conditions in time and space, the usually uneven distribution of the irrigation water, and the seasonal or yearly changes of agricultural practices. Only in lands with undulating topography is the prediction simple: the depressional areas will degrade the most. The preparation of salt and water balances for distinguishable sub-areas in the irrigation project, or the use of agro-hydro-salinity models, can be helpful in explaining or predicting the extent and severity of the problems. Diagnosis Measurement Soil salinity is measured as the salt concentration of the soil solution in tems of g/L or electric conductivity (EC) in dS/m. The relation between these two units is about 5/3: y g/L => 5y/3 dS/m. Seawater may have a salt concentration of 30 g/L (3%) and an EC of 50 dS/m. The standard for the determination of soil salinity is from an extract of a saturated paste of the soil, and the EC is then written as ECe. The extract is obtained by centrifugation. The salinity can more easily be measured, without centrifugation, in a 2:1 or 5:1 water:soil mixture (in terms of g water per g dry soil) than from a saturated paste. The relation between ECe and EC2:1 is about 4, hence: ECe = 4EC1:2. Classification Soils are considered saline when the ECe > 4. When 4 < ECe < 8, the soil is called slightly saline, when 8 < ECe < 16 it is called (moderately) saline, and when ECe > 16 severely saline. Crop tolerance Sensitive crops lose their vigor already in slightly saline soils; most crops are negatively affected by (moderately) saline soils, and only salinity resistant crops thrive in severely saline soils. The University of Wyoming and the Government of Alberta report data on the salt tolerance of plants. Principles of salinity control Drainage is the primary method of controlling soil salinity. The system should permit a small fraction of the irrigation water (about 10 to 20 percent, the drainage or leaching fraction) to be drained and discharged out of the irrigation project. In irrigated areas where salinity is stable, the salt concentration of the drainage water is normally 5 to 10 times higher than that of the irrigation water. Salt export matches salt import and salt will not accumulate. When reclaiming already salinized soils, the salt concentration of the drainage water will initially be much higher than that of the irrigation water (for example 50 times higher). Salt export will greatly exceed salt import, so that with the same drainage fraction a rapid desalinization occurs. After one or two years, the soil salinity is decreased so much, that the salinity of the drainage water has come down to a normal value and a new, favorable, equilibrium is reached. In regions with pronounced dry and wet seasons, the drainage system may be operated in the wet season only, and closed during the dry season. This practice of checked or controlled drainage saves irrigation water. The discharge of salty drainage water may pose environmental problems to downstream areas. The environmental hazards must be considered very carefully and, if necessary mitigating measures must be taken. If possible, the drainage must be limited to wet seasons only, when the salty effluent inflicts the least harm. Drainage systems Land drainage for soil salinity control is usually by horizontal drainage system (figure left), but vertical systems (figure right) are also employed. The drainage system designed to evacuate salty water also lowers the water table. To reduce the cost of the system, the lowering must be reduced to a minimum. The highest permissible level of the water table (or the shallowest permissible depth) depends on the irrigation and agricultural practices and kind of crops. In many cases a seasonal average water table depth of 0.6 to 0.8 m is deep enough. This means that the water table may occasionally be less than 0.6 m (say 0.2 m just after an irrigation or a rain storm). This automatically implies that, in other occasions, the water table will be deeper than 0.8 m (say 1.2 m). The fluctuation of the water table helps in the breathing function of the soil while the expulsion of carbon dioxide (CO2) produced by the plant roots and the inhalation of fresh oxygen (O2) is promoted. The establishing of a not-too-deep water table offers the additional advantage that excessive field irrigation is discouraged, as the crop yield would be negatively affected by the resulting elevated water table, and irrigation water may be saved. The statements made above on the optimum depth of the water table are very general, because in some instances the required water table may be still shallower than indicated (for example in rice paddies), while in other instances it must be considerably deeper (for example in some orchards). The establishment of the optimum depth of the water table is in the realm of agricultural drainage criteria. Soil leaching The vadose zone of the soil below the soil surface and the water table is subject to four main hydrological inflow and outflow factors: Infiltration of rain and irrigation water (Irr) into the soil through the soil surface (Inf) : Inf = Rain + Irr Evaporation of soil water through plants and directly into the air through the soil surface (Evap) Percolation of water from the unsaturated zone soil into the groundwater through the watertable (Perc) Capillary rise of groundwater moving by capillary suction forces into the unsaturated zone (Cap) In steady state (i.e. the amount of water stored in the unsaturated zone does not change in the long run) the water balance of the unsaturated zone reads: Inflow = Outflow, thus: Inf + Cap = Evap + Perc or: Irr + Rain + Cap = Evap + Perc and the salt balance is Irr.Ci + Cap.Cc = Evap.Fc.Ce + Perc.Cp + Ss where Ci is the salt concentration of the irrigation water, Cc is the salt concentration of the capillary rise, equal to the salt concentration of the upper part of the groundwater body, Fc is the fraction of the total evaporation transpired by plants, Ce is the salt concentration of the water taken up by the plant roots, Cp is the salt concentration of the percolation water, and Ss is the increase of salt storage in the unsaturated soil. This assumes that the rainfall contains no salts. Only along the coast this may not be true. Further it is assumed that no runoff or surface drainage occurs. The amount of removed by plants (Evap.Fc.Ce) is usually negligibly small: Evap.Fc.Ce = 0 The salt concentration Cp can be taken as a part of the salt concentration of the soil in the unsaturated zone (Cu) giving: Cp = Le.Cu, where Le is the leaching efficiency. The leaching efficiency is often in the order of 0.7 to 0.8, but in poorly structured, heavy clay soils it may be less. In the Leziria Grande polder in the delta of the Tagus river in Portugal it was found that the leaching efficiency was only 0.15. Assuming that one wishes to avoid the soil salinity to increase and maintain the soil salinity Cu at a desired level Cd we have: Ss = 0, Cu = Cd and Cp = Le.Cd. Hence the salt balance can be simplified to: Perc.Le.Cd = Irr.Ci + Cap.Cc Setting the amount percolation water required to fulfill this salt balance equal to Lr (the leaching requirement) it is found that: Lr = (Irr.Ci + Cap.Cc) / Le.Cd . Substituting herein Irr = Evap + Perc − Rain − Cap and re-arranging gives : Lr = [ (Evap−Rain).Ci + Cap(Cc−Ci) ] / (Le.Cd − Ci) With this the irrigation and drainage requirements for salinity control can be computed too. In irrigation projects in (semi)arid zones and climates it is important to check the leaching requirement, whereby the field irrigation efficiency (indicating the fraction of irrigation water percolating to the underground) is to be taken into account. The desired soil salinity level Cd depends on the crop tolerance to salt. The University of Wyoming, US, and the Government of Alberta, Canada, report crop tolerance data. Strip cropping: an alternative In irrigated lands with scarce water resources suffering from drainage (high water table) and soil salinity problems, strip cropping is sometimes practiced with strips of land where every other strip is irrigated while the strips in between are left permanently fallow. Owing to the water application in the irrigated strips they have a higher water table which induces flow of groundwater to the unirrigated strips. This flow functions as subsurface drainage for the irrigated strips, whereby the water table is maintained at a not-too-shallow depth, leaching of the soil is possible, and the soil salinity can be controlled at an acceptably low level. In the unirrigated (sacrificial) strips the soil is dry and the groundwater comes up by capillary rise and evaporates leaving the salts behind, so that here the soil salinizes. Nevertheless, they can have some use for livestock, sowing salinity resistant grasses or weeds. Moreover, useful salt resistant trees can be planted like Casuarina, Eucalyptus, or Atriplex, keeping in mind that the trees have deep rooting systems and the salinity of the wet subsoil is less than of the topsoil. In these ways wind erosion can be controlled. The unirrigated strips can also be used for salt harvesting. Soil salinity models The majority of the computer models available for water and solute transport in the soil (e.g. SWAP, DrainMod-S, UnSatChem, and Hydrus) are based on Richard's differential equation for the movement of water in unsaturated soil in combination with Fick's differential convection–diffusion equation for advection and dispersion of salts. The models require the input of soil characteristics like the relations between variable unsaturated soil moisture content, water tension, water retention curve, unsaturated hydraulic conductivity, dispersity, and diffusivity. These relations vary greatly from place to place and time to time and are not easy to measure. Further, the models are complicated to calibrate under farmer's field conditions because the soil salinity here is spatially very variable. The models use short time steps and need at least a daily, if not hourly, database of hydrological phenomena. Altogether, this makes model application to a fairly large project the job of a team of specialists with ample facilities. Simpler models, like SaltMod, based on monthly or seasonal water and soil balances and an empirical capillary rise function, are also available. They are useful for long-term salinity predictions in relation to irrigation and drainage practices. LeachMod, Using the SaltMod principles helps in analyzing leaching experiments in which the soil salinity was monitored in various root zone layers while the model will optimize the value of the leaching efficiency of each layer so that a fit is obtained of observed with simulated soil salinity values. Spatial variations owing to variations in topography can be simulated and predicted using salinity cum groundwater models, like SahysMod.	Technology	Soil and soil management	null
5411659	https://en.wikipedia.org/wiki/Marine%20ecosystem	Marine ecosystem	Marine ecosystems are the largest of Earth's aquatic ecosystems and exist in waters that have a high salt content. These systems contrast with freshwater ecosystems, which have a lower salt content. Marine waters cover more than 70% of the surface of the Earth and account for more than 97% of Earth's water supply and 90% of habitable space on Earth. Seawater has an average salinity of 35 parts per thousand of water. Actual salinity varies among different marine ecosystems. Marine ecosystems can be divided into many zones depending upon water depth and shoreline features. The oceanic zone is the vast open part of the ocean where animals such as whales, sharks, and tuna live. The benthic zone consists of substrates below water where many invertebrates live. The intertidal zone is the area between high and low tides. Other near-shore (neritic) zones can include mudflats, seagrass meadows, mangroves, rocky intertidal systems, salt marshes, coral reefs, lagoons. In the deep water, hydrothermal vents may occur where chemosynthetic sulfur bacteria form the base of the food web. Marine ecosystems are characterized by the biological community of organisms that they are associated with and their physical environment. Classes of organisms found in marine ecosystems include brown algae, dinoflagellates, corals, cephalopods, echinoderms, and sharks. Marine ecosystems are important sources of ecosystem services and food and jobs for significant portions of the global population. Human uses of marine ecosystems and pollution in marine ecosystems are significantly threats to the stability of these ecosystems. Environmental problems concerning marine ecosystems include unsustainable exploitation of marine resources (for example overfishing of certain species), marine pollution, climate change, and building on coastal areas. Moreover, much of the carbon dioxide causing global warming and heat captured by global warming are absorbed by the ocean, ocean chemistry is changing through processes like ocean acidification which in turn threatens marine ecosystems. Because of the opportunities in marine ecosystems for humans and the threats created by humans, the international community has prioritized "Life below water" as Sustainable Development Goal 14. The goal is to "Conserve and sustainably use the oceans, seas and marine resources for sustainable development". Types or locations Marine coastal ecosystems Coral reefs Coral reefs are one of the most well-known marine ecosystems in the world, with the largest being the Great Barrier Reef. These reefs are composed of large coral colonies of a variety of species living together. The corals form multiple symbiotic relationships with the organisms around them. Mangroves Mangroves are trees or shrubs that grow in low-oxygen soil near coastlines in tropical or subtropical latitudes. They are an extremely productive and complex ecosystem that connects the land and sea. Mangroves consist of species that are not necessarily related to each other and are often grouped for the characteristics they share rather than genetic similarity. Because of their proximity to the coast, they have all developed adaptions such as salt excretion and root aeration to live in salty, oxygen-depleted water. Mangroves can often be recognized by their dense tangle of roots that act to protect the coast by reducing erosion from storm surges, currents, wave, and tides. The mangrove ecosystem is also an important source of food for many species as well as excellent at sequestering carbon dioxide from the atmosphere with global mangrove carbon storage is estimated at 34 million metric tons per year. Seagrass meadows Seagrasses form dense underwater meadows which are among the most productive ecosystems in the world. They provide habitats and food for a diversity of marine life comparable to coral reefs. This includes invertebrates like shrimp and crabs, cod and flatfish, marine mammals and birds. They provide refuges for endangered species such as seahorses, turtles, and dugongs. They function as nursery habitats for shrimps, scallops and many commercial fish species. Seagrass meadows provide coastal storm protection by the way their leaves absorb energy from waves as they hit the coast. They keep coastal waters healthy by absorbing bacteria and nutrients, and slow the speed of climate change by sequestering carbon dioxide into the sediment of the ocean floor. Seagrasses evolved from marine algae which colonized land and became land plants, and then returned to the ocean about 100 million years ago. However, today seagrass meadows are being damaged by human activities such as pollution from land runoff, fishing boats that drag dredges or trawls across the meadows uprooting the grass, and overfishing which unbalances the ecosystem. Seagrass meadows are currently being destroyed at a rate of about two football fields every hour. Kelp forests Kelp forests occur worldwide throughout temperate and polar coastal oceans. In 2007, kelp forests were also discovered in tropical waters near Ecuador. Physically formed by brown macroalgae, kelp forests provide a unique habitat for marine organisms and are a source for understanding many ecological processes. Over the last century, they have been the focus of extensive research, particularly in trophic ecology, and continue to provoke important ideas that are relevant beyond this unique ecosystem. For example, kelp forests can influence coastal oceanographic patterns and provide many ecosystem services. However, the influence of humans has often contributed to kelp forest degradation. Of particular concern are the effects of overfishing nearshore ecosystems, which can release herbivores from their normal population regulation and result in the overgrazing of kelp and other algae. This can rapidly result in transitions to barren landscapes where relatively few species persist. Already due to the combined effects of overfishing and climate change, kelp forests have all but disappeared in many especially vulnerable places, such as Tasmania's east coast and the coast of Northern California. The implementation of marine protected areas is one management strategy useful for addressing such issues, since it may limit the impacts of fishing and buffer the ecosystem from additive effects of other environmental stressors. Estuaries Estuaries occur where there is a noticeable change in salinity between saltwater and freshwater sources. This is typically found where rivers meet the ocean or sea. The wildlife found within estuaries is unique as the water in these areas is brackish - a mix of freshwater flowing to the ocean and salty seawater. Other types of estuaries also exist and have similar characteristics as traditional brackish estuaries. The Great Lakes are a prime example. There, river water mixes with lake water and creates freshwater estuaries. Estuaries are extremely productive ecosystems that many humans and animal species rely on for various activities. This can be seen as, of the 32 largest cities in the world, 22 are located on estuaries as they provide many environmental and economic benefits such as crucial habitat for many species, and being economic hubs for many coastal communities. Estuaries also provide essential ecosystem services such as water filtration, habitat protection, erosion control, gas regulation nutrient cycling, and it even gives education, recreation and tourism opportunities to people. Lagoons Lagoons are areas that are separated from larger water by natural barriers such as coral reefs or sandbars. There are two types of lagoons, coastal and oceanic/atoll lagoons. A coastal lagoon is, as the definition above, simply a body of water that is separated from the ocean by a barrier. An atoll lagoon is a circular coral reef or several coral islands that surround a lagoon. Atoll lagoons are often much deeper than coastal lagoons. Most lagoons are very shallow meaning that they are greatly affected by changed in precipitation, evaporation and wind. This means that salinity and temperature are widely varied in lagoons and that they can have water that ranges from fresh to hypersaline. Lagoons can be found in on coasts all over the world, on every continent except Antarctica and is an extremely diverse habitat being home to a wide array of species including birds, fish, crabs, plankton and more. Lagoons are also important to the economy as they provide a wide array of ecosystem services in addition to being the home of so many different species. Some of these services include fisheries, nutrient cycling, flood protection, water filtration, and even human tradition. Salt marsh Salt marshes are a transition from the ocean to the land, where fresh and saltwater mix. The soil in these marshes is often made up of mud and a layer of organic material called peat. Peat is characterized as waterlogged and root-filled decomposing plant matter that often causes low oxygen levels (hypoxia). These hypoxic conditions causes growth of the bacteria that also gives salt marshes the sulfurous smell they are often known for. Salt marshes exist around the world and are needed for healthy ecosystems and a healthy economy. They are extremely productive ecosystems and they provide essential services for more than 75 percent of fishery species and protect shorelines from erosion and flooding. Salt marshes can be generally divided into the high marsh, low marsh, and the upland border. The low marsh is closer to the ocean, with it being flooded at nearly every tide except low tide. The high marsh is located between the low marsh and the upland border and it usually only flooded when higher than usual tides are present. The upland border is the freshwater edge of the marsh and is usually located at elevations slightly higher than the high marsh. This region is usually only flooded under extreme weather conditions and experiences much less waterlogged conditions and salt stress than other areas of the marsh. Intertidal zones Intertidal zones are the areas that are visible and exposed to air during low tide and covered up by saltwater during high tide. There are four physical divisions of the intertidal zone with each one having its distinct characteristics and wildlife. These divisions are the Spray zone, High intertidal zone, Middle Intertidal zone, and Low intertidal zone. The Spray zone is a damp area that is usually only reached by the ocean and submerged only under high tides or storms. The high intertidal zone is submerged at high tide but remains dry for long periods between high tides. Due to the large variance of conditions possible in this region, it is inhabited by resilient wildlife that can withstand these changes such as barnacles, marine snails, mussels and hermit crabs. Tides flow over the middle intertidal zone two times a day and this zone has a larger variety of wildlife. The low intertidal zone is submerged nearly all the time except during the lowest tides and life is more abundant here due to the protection that the water gives. Ocean surface Organisms that live freely at the surface, termed neuston, include keystone organisms like the golden seaweed Sargassum that makes up the Sargasso Sea, floating barnacles, marine snails, nudibranchs, and cnidarians. Many ecologically and economically important fish species live as or rely upon neuston. Species at the surface are not distributed uniformly; the ocean's surface harbours unique neustonic communities and ecoregions found at only certain latitudes and only in specific ocean basins. But the surface is also on the front line of climate change and pollution. Life on the ocean's surface connects worlds. From shallow waters to the deep sea, the open ocean to rivers and lakes, numerous terrestrial and marine species depend on the surface ecosystem and the organisms found there. The ocean's surface acts like a skin between the atmosphere above and the water below, and harbours an ecosystem unique to this environment. This sun-drenched habitat can be defined as roughly one metre in depth, as nearly half of UV-B is attenuated within this first meter. Organisms here must contend with wave action and unique chemical and physical properties. The surface is utilised by a wide range of species, from various fish and cetaceans, to species that ride on ocean debris (termed rafters). Most prominently, the surface is home to a unique community of free-living organisms, termed neuston (from the Greek word, υεω, which means both to swim and to float. Floating organisms are also sometimes referred to as pleuston, though neuston is more commonly used). Despite the diversity and importance of the ocean's surface in connecting disparate habitats, and the risks it faces, not a lot is known about neustonic life. A stream of airborne microorganisms circles the planet above weather systems but below commercial air lanes. Some peripatetic microorganisms are swept up from terrestrial dust storms, but most originate from marine microorganisms in sea spray. In 2018, scientists reported that hundreds of millions of viruses and tens of millions of bacteria are deposited daily on every square meter around the planet. Deep sea and sea floor The deep sea contains up to 95% of the space occupied by living organisms. Combined with the sea floor (or benthic zone), these two areas have yet to be fully explored and have their organisms documented. Large marine ecosystems In 1984, National Oceanic and Atmospheric Administration (NOAA) of the United States developed the concept of large marine ecosystems (sometimes abbreviated to LMEs), to identify areas of the oceans for environmental conservation purposes and to enable collaborative ecosystem-based management in transnational areas, in a way consistent with the 1982 UN Convention on the Law of the Sea. This name refers to relatively large regions on the order of or greater, characterized by their distinct bathymetry, hydrography, productivity, and trophically dependent populations. Such LMEs encompass coastal areas from river basins and estuaries to the seaward boundaries of continental shelves and the outer margins of the major ocean current systems. Altogether, there are 66 LMEs, which contribute an estimated $3 trillion annually. This includes being responsible for 90% of global annual marine fishery biomass. LME-based conservation is based on recognition that the world's coastal ocean waters are degraded by unsustainable fishing practices, habitat degradation, eutrophication, toxic pollution, aerosol contamination, and emerging diseases, and that positive actions to mitigate these threats require coordinated actions by governments and civil society to recover depleted fish populations, restore degraded habitats and reduce coastal pollution. Five modules are considered when assessing LMEs: productivity, fish and fisheries, pollution and ecosystem health, socioeconomics, and governance. Periodically assessing the state of each module within a marine LME is encouraged to ensure maintained health of the ecosystem and future benefit to managing governments. The Global Environment Facility (GEF) aids in managing LMEs off the coasts of Africa and Asia by creating resource management agreements between environmental, fisheries, energy and tourism ministers of bordering countries. This means participating countries share knowledge and resources pertaining to local LMEs to promote longevity and recovery of fisheries and other industries dependent upon LMEs. Large marine ecosystems include: East Bering Sea Gulf of Alaska California Current Gulf of California Gulf of Mexico Southeast U.S. Continental Shelf Northeast U.S. Continental Shelf Scotian Shelf Newfoundland-Labrador Shelf Insular Pacific-Hawaiian Pacific Central-American Coastal Caribbean Sea Humboldt Current Patagonian Shelf South Brazil Shelf East Brazil Shelf North Brazil Shelf West Greenland Shelf East Greenland Shelf Barents Sea Norwegian Shelf North Sea Baltic Sea Celtic-Biscay Shelf Central Arctic Iberian Coastal Mediterranean Sea Canary Current Guinea Current Benguela Current Agulhas Current Somali Coastal Current Arabian Sea Red Sea Bay of Bengal Gulf of Thailand South China Sea Sulu-Celebes Sea Indonesian Sea North Australian Shelf Northeast Australian Shelf/Great Barrier Reef East-Central Australian Shelf Southeast Australian Shelf Southwest Australian Shelf West-Central Australian Shelf Northwest Australian Shelf New Zealand Shelf East China Sea Yellow Sea Kuroshio Current Sea of Japan Oyashio Current Sea of Okhotsk West Bering Sea Chukchi Sea Beaufort Sea East Siberian Sea Laptev Sea Kara Sea Iceland Shelf Faroe Plateau Antarctica Black Sea Hudson Bay Arctic Ocean Greenland Sea Role in ecosystem services In addition to providing many benefits to the natural world, marine ecosystems also provide social, economic, and biological ecosystem services to humans. Pelagic marine systems regulate the global climate, contribute to the water cycle, maintain biodiversity, provide food and energy resources, and create opportunities for recreation and tourism. Economically, marine systems support billions of dollars worth of capture fisheries, aquaculture, offshore oil and gas, and trade and shipping. Ecosystem services fall into multiple categories, including supporting services, provisioning services, regulating services, and cultural services. The productivity of a marine ecosystem can be measured in several ways. Measurements pertaining to zooplankton biodiversity and species composition, zooplankton biomass, water-column structure, photosynthetically active radiation, transparency, chlorophyll-a, nitrate, and primary production are used to assess changes in LME productivity and potential fisheries yield. Sensors attached to the bottom of ships or deployed on floats can measure these metrics and be used to quantitatively describe changes in productivity alongside physical changes in the water column such as temperature and salinity. This data can be used in conjunction with satellite measurements of chlorophyll and sea surface temperatures to validate measurements and observe trends on greater spatial and temporal scales. Bottom-trawl surveys and pelagic-species acoustic surveys are used to assess changes in fish biodiversity and abundance in LMEs. Fish populations can be surveyed for stock identification, length, stomach content, age-growth relationships, fecundity, coastal pollution and associated pathological conditions, as well as multispecies trophic relationships. Fish trawls can also collect sediment and inform us about ocean-bottom conditions such as anoxia. Threats Human exploitation and development Coastal marine ecosystems experience growing population pressures with nearly 40% of people in the world living within 100 km of the coast. Humans often aggregate near coastal habitats to take advantage of ecosystem services. For example, coastal capture fisheries from mangroves and coral reef habitats are estimated to be worth a minimum of $34 billion per year. Yet, many of these habitats are either marginally protected or not protected. Mangrove area has declined worldwide by more than one-third since 1950, and 60% of the world's coral reefs are now immediately or directly threatened. Human development, aquaculture, and industrialization often lead to the destruction, replacement, or degradation of coastal habitats. Moving offshore, pelagic marine systems are directly threatened by overfishing. Global fisheries landings peaked in the late 1980s, but are now declining, despite increasing fishing effort. Fish biomass and average trophic level of fisheries landing are decreasing, leading to declines in marine biodiversity. In particular, local extinctions have led to declines in large, long-lived, slow-growing species, and those that have narrow geographic ranges. Biodiversity declines can lead to associated declines in ecosystem services. A long-term study reports the decline of 74–92% of catch per unit effort of sharks in Australian coastline from the 1960s to 2010s. Such biodiversity losses impact not just species themselves, but humans as well, and can contribute to climate change across the globe. The National Oceanic and Atmospheric Administration (NOAA) states that managing and protecting marine ecosystems is crucial in attempting to conserve biodiversity in the face of Earth’s rapidly changing climate. Pollution Invasive species Global aquarium trade Ballast water transport Aquaculture Climate change Warming temperatures (see ocean heat content, sea surface temperature, and marine heat wave) Increased frequency/intensity of storms Ocean acidification Sea level rise Society and culture Global goals By integrating socioeconomic metrics with ecosystem management solutions, scientific findings can be utilized to benefit both the environment and economy of local regions. Management efforts must be practical and cost-effective. In 2000, the Department of Natural Resource Economics at the University of Rhode Island has created a method for measuring and understanding the human dimensions of LMEs and for taking into consideration both socioeconomic and environmental costs and benefits of managing Large Marine Ecosystems. International attention to address the threats of coasts has been captured in Sustainable Development Goal 14 "Life Below Water" which sets goals for international policy focused on preserving coastal ecosystems and supporting more sustainable economic practices for coastal communities. Furthermore, the United Nations has declared 2021-2030 the UN Decade on Ecosystem Restoration, but restoration of coastal ecosystems has received insufficient attention.	Physical sciences	Oceanography	Earth science
2943192	https://en.wikipedia.org/wiki/Coldwater%20fish	Coldwater fish	The term coldwater fish can have different meanings in different contexts. In the context of fishkeeping, it refers to ornamental fish species that tolerate the temperatures of a typical indoor aquarium well and do not require a heater to remain active, as opposed to tropical fish whom need a heater to survive in the room temperatures of temperate climates; In the context of ecology and fishing, it refers to fish species that prefer to inhabit waterbodies or depth zones with much lower temperatures than the average temperate water. Salmonids (e.g. salmon, trout, char and graylings) are a classic example of such types of fish. Fishkeeping Most or all ornamental fish species are able to tolerate temperatures as low as or lower than room temperature, with most stenothermic tropical species having critical thermal minimums of around 10-12 °C. Although these fish are capable of surviving in unheated aquaria, their temperature preferences may vary. For example, koi, goldfish, and pond loaches are commonly considered to be cold-water fish because of their ability to survive at very low temperatures, but their temperature preferences and/or physiological optimal temperatures are , , and , respectively. Because many of the ornamental fish considered to be “coldwater fish” are more accurately eurythermal fish and many prefer temperatures similar to, or even warmer than those preferred by certain tropical fish, the term “coldwater fish” in the aquarium context often misleads pet owners into keeping fish below their preferred temperature. Freshwater aquarium fish Southern redbelly dace Lepomis Shubunkin Comet goldfish Common goldfish Fancy goldfish Black telescope Fantail goldfish Oranda Ryukin Weather loach White Cloud Mountain minnow Celestial Pearl Danio Buenos Aires tetra Gold barb Rosy barb Odessa barb Fathead minnow Banded corydoras Chinese high fin banded shark Three-spined stickleback Ticto barb Pygmy sunfish Enneacanthus Texas cichlid Paradise fish Green barb Zebra danio Bengal danio Leopard danio Danio tinwini Bulldog pleco Rhinogobius Desert goby Highland swordtail (Xiphophorus malinche) Japanese ricefish Zacco Black lined loach (Yasuhikotakia nigrolineata) Red shiner (Cyprinella lutrensis) Spotted gar Longnose gar Rosy red minnow Hillstream loach Spined loach Stone loach Common minnow Vietnamese cardinal minnow GM glowing medaka Gobio Amur bitterling Rosy bitterling Light's bitterling Deep bodied bitterling Rainbow shiner (Notropsis chromosus) Black shark (not to be confused with the tropical red tailed black shark) Golden cobra snakehead Dwarf snakehead Rainbow snakehead Spotted snakehead (Channa punctata) Pearl danio Northern snakehead Chinese algae eater Variable platyfish Note: The above contains a mix of true coldwater fish and sub-tropical fish that can survive and thrive at room temperature which ranges from and to . Freshwater pond fish Three-spined stickleback Nine-spined stickleback Common goldfish Comet goldfish Shubunkin Sterlet Koi Golden orfe Blue orfe Bitterling Gobio Grass carp Albino grass carp Fathead minnow Rosy red minnow Mirror carp Common carp Golden rudd Green tench Golden tench Channel catfish Golden rainbow trout Roach Bluegill Pumpkinseed Weather loach Stone loach Spined loach Common minnow Saltwater aquarium fish Garibaldi Catalina goby Zebra Catalina goby (Lythrypnus zebra) Ornate boxfish Shaw's boxfish White bar boxfish Truncate coralfish Blue devil Pot bellied seahorses Wild fisheries The term "coldwater" is also used to refer to wild fish species that prefer bodies of water that are colder than most temperate waters. In recreational fishing, anglers may loosely break down fish into categories of warm-water fish, cool-water fish, and cold-water fish. Warm-water fish, such as largemouth bass, sunfish and bullhead catfish, are species that tend to dwell in relatively warm tropical and temperate waters similar to the room temperatures that humans easily find comfortable. Cool-water species, such as smallmouth bass and walleye, can tolerate a wide range of temperatures, but tend to be most abundant in cooler rivers or deeper parts of ponds and lakes, where the temperature is slightly lower than room temperatures. Cold-water species, such as salmonids (e.g. salmon, trout, char, graylings, freshwater whitefishes, etc.) and gadiforms (cods, hakes, pollock, haddock, burbot and rocklings, etc.), however become stressed at warm temperatures and are most active in colder temperatures around which resemble a more subarctic or alpine condition. Because these designations are informal, different fisheries management authorities may recognize different boundaries in temperature preference between the categories.	Biology and health sciences	Fishes by habitat	Animals
2943640	https://en.wikipedia.org/wiki/Antibiotic%20sensitivity%20testing	Antibiotic sensitivity testing	Antibiotic sensitivity testing or antibiotic susceptibility testing is the measurement of the susceptibility of bacteria to antibiotics. It is used because bacteria may have resistance to some antibiotics. Sensitivity testing results can allow a clinician to change the choice of antibiotics from empiric therapy, which is when an antibiotic is selected based on clinical suspicion about the site of an infection and common causative bacteria, to directed therapy, in which the choice of antibiotic is based on knowledge of the organism and its sensitivities. Sensitivity testing usually occurs in a medical laboratory, and uses culture methods that expose bacteria to antibiotics, or genetic methods that test to see if bacteria have genes that confer resistance. Culture methods often involve measuring the diameter of areas without bacterial growth, called zones of inhibition, around paper discs containing antibiotics on agar culture dishes that have been evenly inoculated with bacteria. The minimum inhibitory concentration, which is the lowest concentration of the antibiotic that stops the growth of bacteria, can be estimated from the size of the zone of inhibition. Antibiotic susceptibility testing has been needed since the discovery of the beta-lactam antibiotic penicillin. Initial methods were phenotypic, and involved culture or dilution. The Etest, an antibiotic impregnated strip, has been available since the 1980s, and genetic methods such as polymerase chain reaction (PCR) testing have been available since the early 2000s. Research is ongoing into improving current methods by making them faster or more accurate, as well as developing new methods for testing, such as microfluidics. Uses In clinical medicine, antibiotics are most frequently prescribed on the basis of a person's symptoms and medical guidelines. This method of antibiotic selection is called empiric therapy, and it is based on knowledge about what bacteria cause an infection, and to what antibiotics bacteria may be sensitive or resistant. For example, a simple urinary tract infection might be treated with trimethoprim/sulfamethoxazole. This is because Escherichia coli is the most likely causative bacterium, and may be sensitive to that combination antibiotic. However, bacteria can be resistant to several classes of antibiotics. This resistance might be because a type of bacteria has intrinsic resistance to some antibiotics, because of resistance following past exposure to antibiotics, or because resistance may be transmitted from other sources such as plasmids. Antibiotic sensitivity testing provides information about which antibiotics are more likely to be successful and should therefore be used to treat the infection. Antibiotic sensitivity testing is also conducted at a population level in some countries as a form of screening. This is to assess the background rates of resistance to antibiotics (for example with methicillin-resistant Staphylococcus aureus), and may influence guidelines and public health measures. Methods Once a bacterium has been identified following microbiological culture, antibiotics are selected for susceptibility testing. Susceptibility testing methods are based on exposing bacteria to antibiotics and observing the effect on the growth of the bacteria (phenotypic testing), or identifying specific genetic markers (genetic testing). Methods used may be qualitative, meaning that a result indicates resistance is or is not present; or quantitative, using a minimum inhibitory concentration (MIC) to describe the concentration of antibiotic to which a bacterium is sensitive. There are many factors that can affect the results of antibiotic sensitivity testing, including failure of the instrument, temperature, moisture, and potency of the antimicrobial agent. Quality control (QC) testing helps to ensure the accuracy of test results. Organizations such as the American Type Culture Collection and National Collection of Type Cultures provide strains of bacteria with known resistance phenotypes that can be used for quality control. Phenotypic methods Testing based on exposing bacteria to antibiotics uses agar plates or dilution in agar or broth. The selection of antibiotics will depend on the organism grown, and the antibiotics that are available locally. To ensure that the results are accurate, the concentration of bacteria that is added to the agar or broth (the inoculum) must be standardized. This is accomplished by comparing the turbidity of bacteria suspended in saline or broth to McFarland standards—solutions whose turbidity is equivalent to that of a suspension containing a given concentration of bacteria. Once an appropriate concentration (most commonly an 0.5 McFarland standard) has been reached, which can be determined by visual inspection or by photometry, the inoculum is added to the growth medium. Manual The disc diffusion method involves selecting a strain of bacteria, placing it on an agar plate, and observing bacterial growth near antibiotic-impregnated discs. This is also called the Kirby-Bauer method, although modified methods are also used. In some cases, urine samples or positive blood culture samples are applied directly to the test medium, bypassing the preliminary step of isolating the organism. If the antibiotic inhibits microbial growth, a clear ring, or zone of inhibition, is seen around the disc. The bacteria are classified as sensitive, intermediate, or resistant to an antibiotic by comparing the diameter of the zone of inhibition to defined thresholds which correlate with MICs. Mueller–Hinton agar is frequently used in the disc diffusion test. The Clinical and Laboratory Standards Institute (CLSI) and European Committee on Antimicrobial Susceptibility Testing (EUCAST) provide standards for the type and depth of agar, temperature of incubation, and method of analysing results. Disc diffusion is considered the cheapest and most simple of the methods used to test for susceptibility, and is easily adapted to testing newly available antibiotics or formulations. Some slow-growing and fastidious bacteria cannot be accurately tested by this method, while others, such as Streptococcus species and Haemophilus influenzae, can be tested but require specialized growth media and incubation conditions. Gradient methods, such as Etest, use a plastic strip placed on agar. A plastic strip impregnated with different concentrations of antibiotics is placed on a growth medium, and the growth medium is viewed after a period of incubation. The minimum inhibitory concentration can be identified based on the intersection of the teardrop-shaped zone of inhibition with the marking on the strip. Multiple strips for different antibiotics may be used. This type of test is considered a diffusion test. In agar and broth dilution methods, bacteria are placed in multiple small tubes with different concentrations of antibiotics. Whether a bacterium is sensitive or not is determined by visual inspection or automatic optical methods, after a period of incubation. Broth dilution is considered the gold standard for phenotypic testing. The lowest concentration of antibiotics that inhibits growth is considered the MIC. Automated Automated systems exist that replicate manual processes, for example, by using imaging and software analysis to report the zone of inhibition in diffusion testing, or dispensing samples and determining results in dilutional testing. Automated instruments, such as the VITEK 2, BD Phoenix, and Microscan systems, are the most common methodology for AST. The specifications of each instrument vary, but the basic principle involves the introduction of a bacterial suspension into pre-formulated panels of antibiotics. The panels are incubated and the inhibition of bacterial growth by the antibiotic is automatically measured using methodologies such as turbidimetry, spectrophotometry or fluorescence detection. An expert system correlates the MICs with susceptibility results, and the results are automatically transmitted into the laboratory information system for validation and reporting. While such automated testing is less labour-intensive and more standardized than manual testing, its accuracy can be comparatively poor for certain organisms and antibiotics, so the disc diffusion test remains useful as a backup method. Genetic methods Genetic testing, such as via polymerase chain reaction (PCR), DNA microarray, and loop-mediated isothermal amplification, may be used to detect whether bacteria possess genes which confer antibiotic resistance. An example is the use of PCR to detect the mecA gene for beta-lactam resistant Staphylococcus aureus. Other examples include assays for testing vancomycin resistance genes vanA and vanB in Enterococcus species, and antibiotic resistance in Pseudomonas aeruginosa, Klebsiella pneumoniae and Escherichia coli. These tests have the benefit of being direct and rapid, as compared with observable methods, and have a high likelihood of detecting a finding when there is one to detect. However, whether resistance genes are detected does not always match the resistance profile seen with phenotypic method. The tests are also expensive and require specifically trained personnel. Polymerase chain reaction is a method of identifying genes related to antibiotic susceptibility. In the PCR process, a bacterium's DNA is denatured and the two strands of the double helix separate. Primers specific to a sought-after gene are added to a solution containing the DNA, and a DNA polymerase is added alongside a mixture containing molecules that will be needed (for example, nucleotides and ions). If the relevant gene is present, every time this process runs, the quantity of the target gene will be doubled. After this process, the presence of the genes is demonstrated through a variety of methods including electrophoresis, southern blotting, and other DNA sequencing analysis methods. DNA microarrays and chips use the binding of complementary DNA to a target gene or nucleic acid sequence. The benefit of this is that multiple genes can be assessed simultaneously. Using magnetic nanoparticles studded with a beta-2-glycoprotein I peptide imitating a plasma protein, microbial pathogens could selectively be retrieved from blood culture specimens within hours, in a study published September 2024. Magnets are used to fish out the peptide-bacterial complex, followed by genetic testing. MALDI-TOF Matrix-assisted laser desorption ionisation-time of flight mass spectrometry (MALDI-TOF MS) is another method of susceptibility testing. This is a form of time-of-flight mass spectrometry, in which the molecules of a bacterium are subject to matrix-assisted laser desorption. The ionised particles are then accelerated, and spectral peaks recorded, producing an expression profile, which is capable of differentiating specific bacterial strains after being compared to known profiles. This includes, in the context of antibiotic susceptibility testing, strains such as beta-lactamase producing E. coli. MALDI-TOF is rapid and automated. There are limitations to testing in this format however; results may not match the results of phenotypic testing, and acquisition and maintenance is expensive. Reporting Bacteria are marked as sensitive, resistant, or having intermediate resistance to an antibiotic based on the minimum inhibitory concentration (MIC), which is the lowest concentration of the antibiotic that stops the growth of bacteria. The MIC is compared to standard threshold values (called "breakpoints") for a given bacterium and antibiotic. Breakpoints for the same organism and antibiotic may differ based on the site of infection: for example, the CLSI generally defines Streptococcus pneumoniae as sensitive to intravenous penicillin if MICs are ≤0.06 μg/ml, intermediate if MICs are 0.12 to 1 μg/ml, and resistant if MICs are ≥2 μg/ml, but for cases of meningitis, the breakpoints are considerably lower. Sometimes, whether an antibiotic is marked as resistant is also based on bacterial characteristics that are associated with known methods of resistance such as the potential for beta-lactamase production. Specific patterns of drug resistance or multidrug resistance may be noted, such as the presence of an extended-spectrum beta lactamase. Such information may be useful to the clinician, who can change the empiric treatment to a tailored treatment that is directed only at the causative bacterium. The results of antimicrobial susceptibility tests performed during a given time period can be compiled, usually in the form of a table, to form an antibiogram. Antibiograms help the clinician to select the best empiric antimicrobial therapy based on the local resistance patterns until the laboratory test results are available. Clinical practice Ideal antibiotic therapy is based on determining the causal agent and its antibiotic sensitivity. Empiric treatment is often started before laboratory microbiological reports are available. This might be for common or relatively minor infections based on clinical guidelines (such as community-acquired pneumonia), or for serious infections, such as sepsis or bacterial meningitis, in which delayed treatment carries substantial risks. The effectiveness of individual antibiotics varies with the anatomical site of the infection, the ability of the antibiotic to reach the site of infection, and the ability of the bacteria to resist or inactivate the antibiotic. Specimens for antibiotic sensitivity testing are ideally collected before treatment is started. A sample may be taken from the site of a suspected infection; such as a blood culture sample when bacteria are suspected to be present in the bloodstream (bacteraemia), a sputum sample in the case of a pneumonia, or a urine sample in the case of a urinary tract infection. Sometimes multiple samples may be taken if the source of an infection is not clear. These samples are transferred to the microbiology laboratory where they are added to culture media, in or on which the bacteria grow until they are present in sufficient quantities for identification and sensitivity testing to be carried out. When antibiotic sensitivity testing is completed, it will report the organisms present in the sample, and which antibiotics they are susceptible to. Although antibiotic sensitivity testing is done in a laboratory (in vitro), the information provided about this is often clinically relevant to the antibiotics in a person (in vivo). Sometimes, a decision must be made for some bacteria as to whether they are the cause of an infection, or simply commensal bacteria or contaminants, such as Staphylococcus epidermidis and other opportunistic infections. Other considerations may influence the choice of antibiotics, including the need to penetrate through to an infected site (such as an abscess), or the suspicion that one or more causes of an infection were not detected in a sample. History Since the discovery of the beta-lactam antibiotic penicillin, the rates of antimicrobial resistance have increased. Over time, methods for testing the sensitivity of bacteria to antibiotics have developed and changed. Alexander Fleming in the 1920s developed the first method of susceptibility testing. The "gutter method" that he developed was a diffusion method, involving an antibiotic that was diffused through a gutter made of agar. In the 1940s, multiple investigators, including Pope, Foster and Woodruff, Vincent and Vincent used paper discs instead. All these methods involve testing only susceptibility to penicillin. The results were difficult to interpret and not reliable, because of inaccurate results that were not standardised between laboratories. Dilution has been used as a method to grow and identify bacteria since the 1870s, and as a method of testing the susceptibility of bacteria to antibiotics since 1929, also by Alexander Fleming. The way of determining susceptibility changed from how turbid the solution was, to the pH (in 1942), to optical instruments. The use of larger tube-based "macrodilution" testing has been superseded by smaller "microdilution" kits. In 1966, the World Health Organisation confirmed the Kirby–Bauer method as the standard method for susceptibility testing; it is simple, cost-effective and can test multiple antibiotics. The Etest was developed in 1980 by Bolmstrӧm and Eriksson, and MALDI-TOF developed in 2000s. An array of automated systems has been developed since and after the 1980s. PCR was the first genetic test available and first published as a method of detecting antibiotic susceptibility in 2001. Further research Point-of-care testing is being developed to speed up the time for testing, and to help practitioners avoid prescribing unnecessary antibiotics in the style of precision medicine. Traditional techniques typically take between 12 and 48 hours, although it can take up to five days. In contrast, rapid testing using molecular diagnostics is defined as "being feasible within an 8-h(our) working shift". Progress has been slow due to a range of reasons including cost and regulation. Additional research is focused at the shortcomings of current testing methods. As well as the duration it takes to report phenotypic methods, they are laborious, have difficult portability and are difficult to use in resource-limited settings, and have a chance of cross-contamination. As of 2017, point-of-care resistance diagnostics were available for methicillin-resistant Staphylococcus aureus (MRSA), rifampin-resistant Mycobacterium tuberculosis (TB), and vancomycin-resistant enterococci (VRE) through GeneXpert by molecular diagnostics company Cepheid. Quantitative PCR, with the view of determining the percent of a detected bacteria that possesses a resistance gene, is being explored. Whole genome sequencing of isolated bacteria is also being explored, and likely to become more available as costs decrease and speed increases over time. Additional methods explored include microfluidics, which uses a small amount of fluid and a variety of testing methods, such as optical, electrochemical, and magnetic. Such assays do not require much fluid to be tested, are rapid and portable. The use of fluorescent dyes has been explored. These involve labelled proteins targeted at biomarkers, nucleic acid sequences present within cells that are found when the bacterium is resistant to an antibiotic. An isolate of bacteria is fixed in position and then dissolved. The isolate is then exposed to fluorescent dye, which will be luminescent when viewed. Improvements to existing platforms are also being explored, including improvements in imaging systems that are able to more rapidly identify the MIC in phenotypic samples; or the use of bioluminescent enzymes that reveal bacterial growth to make changes more easily visible. Bibliography	Biology and health sciences	Basics_3	Biology
311410	https://en.wikipedia.org/wiki/Pregnancy%20test	Pregnancy test	A pregnancy test is used to determine whether a female is pregnant or not. The two primary methods are testing for the female pregnancy hormone (human chorionic gonadotropin (hCG)) in blood or urine using a pregnancy test kit, and scanning with ultrasonography. Testing blood for hCG results in the earliest detection of pregnancy. Almost all pregnant women will have a positive urine pregnancy test one week after the first day of a missed menstrual period. Types Human chorionic gonadotropin (hCG) Identified in the early 20th century, human chorionic gonadotropin (hCG) is a glycoprotein hormone that rises quickly in the first few weeks of pregnancy, typically reaching a peak at 8- to 10-weeks gestational age. hCG is produced by what will become the placenta. hCG testing can be performed with a blood (serum) sample (typically done in a medical facility) or with urine (which can be performed in a medical facility or at home). The assays used to detect the presence of hCG in blood or urine are generally reliable and inexpensive. Secretion of hCG can occur as soon as 6 days following ovulation and on average 8–10 days following ovulation; this is the earliest hCG can be detected in a blood sample. The hCG concentration in blood is higher than in urine. Therefore, a blood test can be positive while the urine test is still negative. Qualitative tests (yes/no or positive/negative results) look for the presence of the beta subunit of human chorionic gonadotropin in blood or urine. For a qualitative test the thresholds for a positive test are generally determined by an hCG cut-off where at least 95% of pregnant women would get a positive result on the day of their first missed period. Qualitative urine pregnancy tests vary in sensitivity. High-sensitivity tests are more common and typically detect hCG levels between 20 and 50 milli-international units/mL (mIU/mL). Low-sensitivity tests detect hCG levels between 1500 and 2000 mIU/mL and have unique clinical applications, including confirmation of medication abortion success. Qualitative urine tests available for home use are typically designed as lateral flow tests. Quantitative tests measure the exact amount of hCG in the sample. Blood tests can detect hCG levels as low as 1 mIU/mL, and typically clinicians will diagnose a positive pregnancy test at 5mIU/mL. There is a multilevel urine pregnancy test (MLPT) that measures hCG levels semiquantitatively. The hCG levels are measured at <25, 25 to 99, 100 to 499, 500 to 1999, 2000 to 9999, and >10,000 mIU/mL. This test has utility for determining the success of medication abortion. Ultrasound Obstetric ultrasonography may also be used to detect and diagnose pregnancy. It is very common to have a positive at-home urine pregnancy test before an ultrasound. Both abdominal and vaginal ultrasound may be used, but vaginal ultrasound allows for earlier visualization of the pregnancy. With obstetric ultrasonography the gestational sac (intrauterine fluid collection) can be visualized at 4.5 to 5 weeks gestation, the yolk sac at 5 to 6 weeks gestation, and fetal pole at 5.5 to 6 weeks gestation. Ultrasound is used to diagnose multiple gestation, which cannot be diagnosed based on the presence of hCG in urine or blood. Determination of the gestational age of the embryo/fetus is an additional benefit of ultrasound compared to hCG tests. Accuracy A systematic review published in 1998 showed that home pregnancy test kits, when used by experienced technicians, are almost as accurate as professional laboratory testing (97.4%). When used by consumers, however, the accuracy fell to 75%: the review authors noted that many users misunderstood or failed to follow the instructions included in the kits. False positive False positive pregnancy test results are rare and may occur for several reasons, including: user error in performing and interpreting the test, biochemical pregnancy (loss of pregnancy before signs of pregnancy are apparent on ultrasound, likely very soon after implantation), and non-pregnant production of the hCG molecule (i.e. secretion due to a tumor or the pituitary gland, some diseases of the liver, cancers, including choriocarcinoma and other germ cell tumors, IgA deficiencies, heterophile antibodies, enterocystoplasties, gestational trophoblastic diseases (GTD), and gestational trophoblastic neoplasms). bacterial contamination and blood in urine Spurious evaporation lines may appear on many home pregnancy tests if read after the suggested 3–5 minute window or reaction time, independent of an actual pregnancy. False positives may also appear on tests used past their expiration date. False positive pregnancy test can happen due to 'phantom hCG' which is due to people having human antianimal or heterophilic antibodies. False positives can also be caused by (in order of incidence) quiescent pregnancy, pituitary sulfated hCG, heterophilic antibody, familial hCG syndrome and cancer. Due to use of medication Urine tests can be falsely positive in those that are taking the medications: chlorpromazine, promethazine, phenothiazines, methadone, aspirin, carbamazepine and drugs that cause high urinary pH. False negative False negative readings can occur when testing is done too early. hCG levels rise rapidly in early pregnancy and the chances of false negative test results diminish with time (increasing gestational age). Less sensitive urine tests and qualitative blood tests may not detect pregnancy until three or four days after implantation. Menstruation occurs on average 14 days after ovulation, so the likelihood of a false negative is low once a menstrual period is late. Ovulation may not occur at a predictable time in the menstrual cycle. A number of factors may cause an unexpectedly early or late ovulation, even for people with a history of regular menstrual cycles. Medical providers often struggle to 'rule out' pregnancy for medical testing or treatment that cannot be conducted during pregnancy before they can do an accurate urine pregnancy test. More rare, false negative results can also occur due to a "hook effect", where a sample with a very high level of hCG is tested without dilution, causing an invalid result. Other uses Pregnancy tests may be used to predict if a pregnancy is likely to continue or is abnormal. Miscarriage, or spontaneous abortion or pregnancy loss, is common in early pregnancy. Serial quantitative blood tests may be done, usually 48 hours apart, and interpreted based on the knowledge that hCG in a viable normal pregnancy rises rapidly in early pregnancy. For example, for a starting hCG level of 1,500 mIU/ml or less, the hCG of continuing, normal pregnancy will increase at least 49% in 48 hours. However, for pregnancies with a higher starting hCG, between 1,500 and 3,000 mIU/ml, the hCG should rise at least 40%; for a starting hCG greater than 3,000 mIU/ml, the hCG should increase at least 33%. Failure to rise by these minimums may indicate that the pregnancy is not normal, either as a failed intrauterine pregnancy or a possible ectopic pregnancy. Ultrasound is also a common tool for determining viability and location of a pregnancy. Serial ultrasound may be used to identify non-viable pregnancies, as pregnancies that do not grow in size or develop expected structural findings on repeated ultrasounds over a 1–2 week interval may be identified as abnormal. Occasionally, a single ultrasound may be used to identify a pregnancy as non-viable; for example, an embryo that is greater than a certain size but that lacks a visible heart beat may be confidently determined to be not viable without the need for follow up ultrasound for confirmation. Research Research has identified at least one other possible marker that may appear earlier and exclusively during pregnancy. For example, early pregnancy factor (EPF) can be detected in blood within 48 hours of fertilization, rather than after implantation. However, its reliable use as a pregnancy test remains unclear as studies have shown its presence in physiological situations besides pregnancy, and its application to humans remains limited. History Records of attempts at pregnancy testing have been found as far back as the ancient Greek and ancient Egyptian cultures. The ancient Egyptians watered bags of wheat and barley with the urine of a possibly pregnant woman. Germination indicated pregnancy. The type of grain that sprouted was taken as an indicator of the fetus's sex. Hippocrates suggested that a woman who had missed her period should drink a solution of honey in water at bedtime: resulting abdominal distention and cramps would indicate the presence of a pregnancy. Avicenna and many physicians after him in the Middle Ages performed uroscopy, a nonscientific method to evaluate urine. Selmar Aschheim and Bernhard Zondek introduced testing based on the presence of human chorionic gonadotropin (hCG) in 1928. Early studies of hCG had concluded that it was produced by the pituitary gland. In the 1930s, Doctor Georgeanna Jones discovered that hCG was produced not by the pituitary gland, but by the placenta. This discovery was important in relying on hCG as an early marker of pregnancy. In the Aschheim and Zondek test, an infantile female mouse was injected subcutaneously with urine of the woman to be tested, and the mouse later was killed and dissected. Presence of ovulation indicated that the urine contained hCG and meant that the subject was pregnant. A similar test was developed using immature rabbits. Here, too, killing the animal to check her ovaries was necessary. At the beginning of the 1930s, Hillel Shapiro and Harry Zwarenstein, who were researchers at the University of Cape Town, discovered that if urine from a pregnant woman was injected into the South African Xenopus frog and the frog ovulated, this indicated that the subject was pregnant. This test, known as the frog test, was used throughout the world from the 1930s to 1960s, with Xenopus frogs being exported in great numbers. Shapiro's advisor, Lancelot Hogben, claimed to have developed the pregnancy test himself, but this was refuted by both Shapiro and Zwarenstein in a letter to the British Medical Journal. A later article, independently authored, granted Hogben credit for the principle of using Xenopus to determine gonadotropin levels in a pregnant woman's urine, but not for its usage as a functional pregnancy test. Hormonal pregnancy tests such as Primodos and Duogynon were used in the 1960s and 1970s in the UK and Germany. These tests involved taking a dosed amount of hormones, and observing the response a few days later. A pregnant woman does not react, as she is producing the hormones in pregnancy; a subject who is not pregnant responds to the absence of the hormone by beginning a new menstrual cycle. While the test was (is) generally considered accurate, research advancements have replaced it with simpler techniques. Immunologic pregnancy tests were introduced in 1960 when Wide and Gemzell presented a test based on in-vitro hemagglutination inhibition. This was a first step away from in-vivo pregnancy testing and initiated a series of improvements in pregnancy testing leading to the contemporary at-home testing. Direct measurement of antigens, such as hCG, was made possible after the invention of the radioimmunoassay in 1959. Radioimmunoassays require sophisticated apparatus and special radiation precautions and are expensive. Organon International obtained the first patent on a home pregnancy test in 1969, two years after product designer Margaret Crane noticed that the laboratory testing procedure was relatively simple and made a prototype. The product became available in Canada in 1971, and the United States in 1977, after delays caused by concerns over sexual morality and the ability of potentially pregnant women to perform the test and cope with the results without a doctor. Another home pregnancy testing kit was based on the work of Judith Vaitukaitis and Glenn Braunstein, who developed a sensitive hCG assay at the National Institutes of Health. That test went onto the market in 1978. In the 1970s, the discovery of monoclonal antibodies led to the development of the relatively simple and cheap immunoassays, such as agglutination-inhibition-based assays and sandwich ELISA, used in modern home pregnancy tests. Tests are now so cheap that they can be mass-produced in a general publication and used for advertising.	Biology and health sciences	Diagnostics	Health
311440	https://en.wikipedia.org/wiki/Mammary%20gland	Mammary gland	A mammary gland is an exocrine gland in humans and other mammals that produces milk to feed young offspring. Mammals get their name from the Latin word mamma, "breast". The mammary glands are arranged in organs such as the breasts in primates (for example, humans and chimpanzees), the udder in ruminants (for example, cows, goats, sheep, and deer), and the dugs of other animals (for example, dogs and cats). Lactorrhea, the occasional production of milk by the glands, can occur in any mammal, but in most mammals, lactation, the production of enough milk for nursing, occurs only in phenotypic females who have gestated in recent months or years. It is directed by hormonal guidance from sex steroids. In a few mammalian species, male lactation can occur. With humans, male lactation can occur only under specific circumstances. Mammals are divided into 3 groups: prototherians, metatherians, and eutherians. In the case of prototherians, both males and females have functional mammary glands, but their mammary glands are without nipples. These mammary glands are modified sebaceous glands. Concerning most metatherians and eutherians, only females have functional mammary glands, with the exception of some bat species. Their mammary glands can be termed as breasts or udders. In the case of breasts, each mammary gland has its own nipple (e.g., human mammary glands). In the case of udders, pairs of mammary glands comprise a single mass, with more than one nipple (or teat) hanging from it. For instance, cows and buffalo udders have two pairs of mammary glands and four teats, whereas sheep and goat udders have one pair of mammary glands with two teats protruding from the udder. Each gland produces milk for a single teat. These mammary glands are evolutionarily derived from sweat glands. Structure The basic components of a mature mammary gland are the alveoli (hollow cavities, a few millimeters large), which are lined with milk-secreting cuboidal cells and surrounded by myoepithelial cells. These alveoli join to form groups known as lobules. Each lobule has a lactiferous duct that drains into openings in the nipple. The myoepithelial cells contract under the stimulation of oxytocin, excreting the milk secreted by alveolar units into the lobule lumen toward the nipple. As the infant begins to suck, the oxytocin-mediated "let down reflex" ensues, and the mother's milk is secreted—not sucked—from the gland into the infant's mouth. All the milk-secreting tissue leading to a single lactiferous duct is collectively called a "simple mammary gland"; in a "complex mammary gland", all the simple mammary glands serve one nipple. Humans normally have two complex mammary glands, one in each breast, and each complex mammary gland consists of 10–20 simple glands. The opening of each simple gland on the surface of the nipple is called a "pore." The presence of more than two nipples is known as polythelia and the presence of more than two complex mammary glands as polymastia. Maintaining the correct polarized morphology of the lactiferous duct tree requires another essential component – mammary epithelial cells extracellular matrix (ECM) which, together with adipocytes, fibroblast, inflammatory cells, and others, constitute mammary stroma. Mammary epithelial ECM mainly contains myoepithelial basement membrane and the connective tissue. They not only help to support mammary basic structure, but also serve as a communicating bridge between mammary epithelia and their local and global environment throughout this organ's development. Histology A mammary gland is a specific type of apocrine gland specialized for manufacture of colostrum (first milk) when giving birth. Mammary glands can be identified as apocrine because they exhibit striking "decapitation" secretion. Many sources assert that mammary glands are modified sweat glands. Development Mammary glands develop during different growth cycles. They exist in both sexes during the embryonic stage, forming only a rudimentary duct tree at birth. In this stage, mammary gland development depends on systemic (and maternal) hormones, but is also under the (local) regulation of paracrine communication between neighboring epithelial and mesenchymal cells by parathyroid hormone-related protein (PTHrP). This locally secreted factor gives rise to a series of outside-in and inside-out positive feedback between these two types of cells, so that mammary bud epithelial cells can proliferate and sprout down into the mesenchymal layer until they reach the fat pad to begin the first round of branching. At the same time, the embryonic mesenchymal cells around the epithelial bud receive secreting factors activated by PTHrP, such as BMP4. These mesenchymal cells can transform into a dense, mammary-specific mesenchyme, which later develop into connective tissue with fibrous threads, forming blood vessels and the lymph system. A basement membrane, mainly containing laminin and collagen, formed afterward by differentiated myoepithelial cells, keeps the polarity of this primary duct tree. These components of the extracellular matrix are strong determinants of duct morphogenesis. Biochemistry Estrogen and growth hormone (GH) are essential for the ductal component of mammary gland development, and act synergistically to mediate it. Neither estrogen nor GH are capable of inducing ductal development without the other. The role of GH in ductal development has been found to be mostly mediated by its induction of the secretion of insulin-like growth factor 1 (IGF-1), which occurs both systemically (mainly originating from the liver) and locally in the mammary fat pad through activation of the growth hormone receptor (GHR). However, GH itself also acts independently of IGF-1 to stimulate ductal development by upregulating estrogen receptor (ER) expression in mammary gland tissue, which is a downstream effect of mammary gland GHR activation. In any case, unlike IGF-1, GH itself is not essential for mammary gland development, and IGF-1 in conjunction with estrogen can induce normal mammary gland development without the presence of GH. In addition to IGF-1, other paracrine growth factors such as epidermal growth factor (EGF), transforming growth factor beta (TGF-β), amphiregulin, fibroblast growth factor (FGF), and hepatocyte growth factor (HGF) are involved in breast development as mediators downstream to sex hormones and GH/IGF-1. During embryonic development, IGF-1 levels are low, and gradually increase from birth to puberty. At puberty, the levels of GH and IGF-1 reach their highest levels in life and estrogen begins to be secreted in high amounts in females, which is when ductal development mostly takes place. Under the influence of estrogen, stromal and fat tissue surrounding the ductal system in the mammary glands also grows. After puberty, GH and IGF-1 levels progressively decrease, which limits further development until pregnancy, if it occurs. During pregnancy, progesterone and prolactin are essential for mediating lobuloalveolar development in estrogen-primed mammary gland tissue, which occurs in preparation of lactation and nursing. Androgens such as testosterone inhibit estrogen-mediated mammary gland development (e.g., by reducing local ER expression) through activation of androgen receptors expressed in mammary gland tissue, and in conjunction with relatively low estrogen levels, are the cause of the lack of developed mammary glands in males. Timeline Before birth Mammary gland development is characterized by the unique process by which the epithelium invades the stroma. The development of the mammary gland occurs mainly after birth. During puberty, tubule formation is coupled with branching morphogenesis which establishes the basic arboreal network of ducts emanating from the nipple. Developmentally, mammary gland epithelium is constantly produced and maintained by rare epithelial cells, dubbed as mammary progenitors which are ultimately thought to be derived from tissue-resident stem cells. Embryonic mammary gland development can be divided into a series of specific stages. Initially, the formation of the milk lines that run between the fore and hind limbs bilaterally on each side of the midline occurs around embryonic day 10.5 (E10.5). The second stage occurs at E11.5 when placode formation begins along the mammary milk line. This will eventually give rise to the nipple. Lastly, the third stage occurs at E12.5 and involves the invagination of cells within the placode into the mesenchyme, leading to a mammary anlage (biology). The primitive (stem) cells are detected in embryo and their numbers increase steadily during development Growth Postnatally, the mammary ducts elongate into the mammary fat pad. Then, starting around four weeks of age, mammary ductal growth increases significantly with the ducts invading towards the lymph node. Terminal end buds, the highly proliferative structures found at the tips of the invading ducts, expand and increase greatly during this stage. This developmental period is characterized by the emergence of the terminal end buds and lasts until an age of about 7–8 weeks. By the pubertal stage, the mammary ducts have invaded to the end of the mammary fat pad. At this point, the terminal end buds become less proliferative and decrease in size. Side branches form from the primary ducts and begin to fill the mammary fat pad. Ductal development decreases with the arrival of sexual maturity and undergoes estrous cycles (proestrus, estrus, metestrus, and diestrus). As a result of estrous cycling, the mammary gland undergoes dynamic changes where cells proliferate and then regress in an ordered fashion. Pregnancy During pregnancy, the ductal systems undergo rapid proliferation and form alveolar structures within the branches to be used for milk production. After delivery, lactation occurs within the mammary gland; lactation involves the secretion of milk by the luminal cells in the alveoli. Contraction of the myoepithelial cells surrounding the alveoli will cause the milk to be ejected through the ducts and into the nipple for the nursing infant. Upon weaning of the infant, lactation stops and the mammary gland turns in on itself, a process called involution. This process involves the controlled collapse of mammary epithelial cells where cells begin apoptosis in a controlled manner, reverting the mammary gland back to a pubertal state. Postmenopausal During postmenopause, due to much lower levels of estrogen, and due to lower levels of GH and IGF-1, which decrease with age, mammary gland tissue atrophies and the mammary glands become smaller. Physiology Hormonal control Lactiferous duct development occurs in females in response to circulating hormones. First development is frequently seen during pre- and postnatal stages, and later during puberty. Estrogen promotes branching differentiation, whereas in males testosterone inhibits it. A mature duct tree reaching the limit of the fat pad of the mammary gland comes into being by bifurcation of duct terminal end buds (TEB), secondary branches sprouting from primary ducts and proper duct lumen formation. These processes are tightly modulated by components of mammary epithelial ECM interacting with systemic hormones and local secreting factors. However, for each mechanism the epithelial cells' "niche" can be delicately unique with different membrane receptor profiles and basement membrane thickness from specific branching area to area, so as to regulate cell growth or differentiation sub-locally. Important players include beta-1 integrin, epidermal growth factor receptor (EGFR), laminin-1/5, collagen-IV, matrix metalloproteinase (MMPs), heparan sulfate proteoglycans, and others. Elevated circulating level of growth hormone and estrogen get to multipotent cap cells on TEB tips through a thin, leaky layer of basement membrane. These hormones promote specific gene expression. Hence cap cells can differentiate into myoepithelial and luminal (duct) epithelial cells, and the increased amount of activated MMPs can degrade surrounding ECM helping duct buds to reach further in the fat pads. On the other hand, basement membrane along the mature mammary ducts is thicker, with strong adhesion to epithelial cells via binding to integrin and non-integrin receptors. When side branches develop, it is a much more "pushing-forward" working process including extending through myoepithelial cells, degrading basement membrane and then invading into a periductal layer of fibrous stromal tissue. Degraded basement membrane fragments (laminin-5) roles to lead the way of mammary epithelial cells migration. Whereas, laminin-1 interacts with non-integrin receptor dystroglycan negatively regulates this side branching process in case of cancer. These complex "Yin-yang" balancing crosstalks between mammary ECM and epithelial cells "instruct" healthy mammary gland development until adult. There is preliminary evidence that soybean intake mildly stimulates the breast glands in pre- and postmenopausal women. Pregnancy Secretory alveoli develop mainly in pregnancy, when rising levels of prolactin, estrogen, and progesterone cause further branching, together with an increase in adipose tissue and a richer blood flow. In gestation, serum progesterone remains at a stably high concentration so signaling through its receptor is continuously activated. As one of the transcribed genes, Wnts secreted from mammary epithelial cells act paracrinely to induce more neighboring cells' branching. When the lactiferous duct tree is almost ready, "leaves" alveoli are differentiated from luminal epithelial cells and added at the end of each branch. In late pregnancy and for the first few days after giving birth, colostrum is secreted. Milk secretion (lactation) begins a few days later due to reduction in circulating progesterone and the presence of another important hormone prolactin, which mediates further alveologenesis, milk protein production, and regulates osmotic balance and tight junction function. Laminin and collagen in myoepithelial basement membrane interacting with beta-1 integrin on epithelial surface again, is essential in this process. Their binding ensures correct placement of prolactin receptors on the basal lateral side of alveoli cells and directional secretion of milk into lactiferous ducts. Suckling of the baby causes release of the hormone oxytocin, which stimulates contraction of the myoepithelial cells. In this combined control from ECM and systemic hormones, milk secretion can be reciprocally amplified so as to provide enough nutrition for the baby. Weaning During weaning, decreased prolactin, missing mechanical stimulation (baby suckling), and changes in osmotic balance caused by milk stasis and leaking of tight junctions cause cessation of milk production. It is the (passive) process of a child or animal ceasing to be dependent on the mother for nourishment. In some species there is complete or partial involution of alveolar structures after weaning, in humans there is only partial involution and the level of involution in humans appears to be highly individual. The glands in the breast do secrete fluid also in nonlactating women. In some other species (such as cows), all alveoli and secretory duct structures collapse by programmed cell death (apoptosis) and autophagy for lack of growth promoting factors either from the ECM or circulating hormones. At the same time, apoptosis of blood capillary endothelial cells speeds up the regression of lactation ductal beds. Shrinkage of the mammary duct tree and ECM remodeling by various proteinase is under the control of somatostatin and other growth inhibiting hormones and local factors. This major structural change leads loose fat tissue to fill the empty space afterward. But a functional lactiferous duct tree can be formed again when a female is pregnant again. Clinical significance Tumorigenesis in mammary glands can be induced biochemically by abnormal expression level of circulating hormones or local ECM components, or from a mechanical change in the tension of mammary stroma. Under either of the two circumstances, mammary epithelial cells would grow out of control and eventually result in cancer. Almost all instances of breast cancer originate in the lobules or ducts of the mammary glands. Other mammals General The breasts of female humans vary from most other mammals that tend to have less conspicuous mammary glands. The number and positioning of mammary glands varies widely in different mammals. The protruding teats and accompanying glands can be located anywhere along the two milk lines. In general most mammals develop mammary glands in pairs along these lines, with a number approximating the number of young typically birthed at a time. The number of teats varies from 2 (in most primates) to 18 (in pigs). The Virginia opossum has 13, one of the few mammals with an odd number. The following table lists the number and position of teats and glands found in a range of mammals: Male mammals typically have rudimentary mammary glands and nipples, with a few exceptions: male mice do not have nipples, male marsupials do not have mammary glands, and male horses lack nipples. The male dayak fruit bat has lactating mammary glands. Male lactation occurs infrequently in some species. Mammary glands are true protein factories, and several labs have constructed transgenic animals, mainly goats and cows, to produce proteins for pharmaceutical use. Complex glycoproteins such as monoclonal antibodies or antithrombin cannot be produced by genetically engineered bacteria, and the production in live mammals is much cheaper than the use of mammalian cell cultures. Evolution There are many theories on how mammary glands evolved. For example, it is thought that the mammary gland is a transformed sweat gland, more closely related to apocrine sweat glands. Because mammary glands do not fossilize well, supporting such theories with fossil evidence is difficult. Many of the current theories are based on comparisons between lines of living mammals—monotremes, marsupials, and eutherians. One theory proposes that mammary glands evolved from glands that were used to keep the eggs of early mammals moist and free from infection (monotremes still lay eggs). Other theories suggest that early secretions were used directly by hatched young, or that the secretions were used by young to help them orient to their mothers. Lactation is thought to have developed long before the evolution of the mammary gland and mammals; see evolution of lactation. Additional images	Biology and health sciences	Integumentary system	Biology
311509	https://en.wikipedia.org/wiki/Bounded%20function	Bounded function	In mathematics, a function defined on some set with real or complex values is called bounded if the set of its values is bounded. In other words, there exists a real number such that for all in . A function that is not bounded is said to be unbounded. If is real-valued and for all in , then the function is said to be bounded (from) above by . If for all in , then the function is said to be bounded (from) below by . A real-valued function is bounded if and only if it is bounded from above and below. An important special case is a bounded sequence, where is taken to be the set of natural numbers. Thus a sequence is bounded if there exists a real number such that for every natural number . The set of all bounded sequences forms the sequence space . The definition of boundedness can be generalized to functions taking values in a more general space by requiring that the image is a bounded set in . Related notions Weaker than boundedness is local boundedness. A family of bounded functions may be uniformly bounded. A bounded operator is not a bounded function in the sense of this page's definition (unless ), but has the weaker property of preserving boundedness; bounded sets are mapped to bounded sets . This definition can be extended to any function if and allow for the concept of a bounded set. Boundedness can also be determined by looking at a graph. Examples The sine function is bounded since for all . The function , defined for all real except for −1 and 1, is unbounded. As approaches −1 or 1, the values of this function get larger in magnitude. This function can be made bounded if one restricts its domain to be, for example, or . The function , defined for all real , is bounded, since for all . The inverse trigonometric function arctangent defined as: or is increasing for all real numbers and bounded with radians By the boundedness theorem, every continuous function on a closed interval, such as , is bounded. More generally, any continuous function from a compact space into a metric space is bounded. All complex-valued functions which are entire are either unbounded or constant as a consequence of Liouville's theorem. In particular, the complex must be unbounded since it is entire. The function which takes the value 0 for rational number and 1 for irrational number (cf. Dirichlet function) is bounded. Thus, a function does not need to be "nice" in order to be bounded. The set of all bounded functions defined on is much larger than the set of continuous functions on that interval. Moreover, continuous functions need not be bounded; for example, the functions and defined by and are both continuous, but neither is bounded. (However, a continuous function must be bounded if its domain is both closed and bounded.)	Mathematics	Functions: General	null
311596	https://en.wikipedia.org/wiki/Allicin	Allicin	Allicin is an organosulfur compound obtained from garlic and leeks. When fresh garlic is chopped or crushed, the enzyme alliinase converts alliin into allicin, which is responsible for the aroma of fresh garlic. Allicin is unstable and quickly changes into a series of other sulfur-containing compounds such as diallyl disulfide. Allicin is an antifeedant, i.e. the defense mechanism against attacks by pests on the garlic plant. Allicin is an oily, slightly yellow liquid that gives garlic its distinctive odor. It is a thioester of sulfenic acid. It is also known as allyl thiosulfinate. Its biological activity can be attributed to both its antioxidant activity and its reaction with thiol-containing proteins. Structure and occurrence Allicin features the thiosulfinate functional group, R-S-(O)-S-R. The compound is not present in garlic unless tissue damage occurs, and is formed by the action of the enzyme alliinase on alliin. Allicin is chiral but occurs naturally only as a racemate. The racemic form can also be generated by oxidation of diallyl disulfide: (SCH2CH=CH2)2 + 2 RCO3H + H2O → 2 CH2=CHCH2SOH + 2 RCO2H 2 CH2=CHCH2SOH → CH2=CHCH2S(O)SCH2CH=CH2 + H2O Alliinase is irreversibly deactivated below pH 3; as such, allicin is generally not produced in the body from the consumption of fresh or powdered garlic. Furthermore, allicin can be unstable, breaking down within 16 hours at 23 °C. Biosynthesis The biosynthesis of allicin commences with the conversion of cysteine into S-allyl-L-cysteine. Oxidation of this thioether gives the sulfoxide (alliin). The enzyme alliinase, which contains pyridoxal phosphate (PLP), cleaves alliin, generating allylsulfenic acid (CH2=CHCH2SOH), pyruvate, and ammonium ions. At room temperature, two molecules of allylsulfenic acid condense to form allicin. Research Allicin has been studied for its potential to treat various kinds of multiple drug resistance bacterial infections, as well as viral and fungal infections in vitro, but as of 2016, the safety and efficacy of allicin to treat infections in people was unclear. A Cochrane review found there to be insufficient clinical evidence regarding the effects of allicin in preventing or treating common cold. History It was first isolated and studied in the laboratory by Chester J. Cavallito and John Hays Bailey in 1944. Allicin was discovered as part of efforts to create thiamine derivatives in the 1940s, mainly in Japan. Allicin became a model for medicinal chemistry efforts to create other thiamine disulfides. The results included sulbutiamine, fursultiamine (thiamine tetrahydrofurfuryl disulfide) and benfothiamine. These compounds are hydrophobic, easily pass from the intestines to the bloodstream, and are reduced to thiamine by cysteine or glutathione.	Physical sciences	Concepts: General	Chemistry
311888	https://en.wikipedia.org/wiki/Cornea	Cornea	The cornea is the transparent front part of the eye that covers the iris, pupil, and anterior chamber. Along with the anterior chamber and lens, the cornea refracts light, accounting for approximately two-thirds of the eye's total optical power. In humans, the refractive power of the cornea is approximately 43 dioptres. The cornea can be reshaped by surgical procedures such as LASIK. While the cornea contributes most of the eye's focusing power, its focus is fixed. Accommodation (the refocusing of light to better view near objects) is accomplished by changing the geometry of the lens. Medical terms related to the cornea often start with the prefix "kerat-" from the Greek word κέρας, horn. Structure The cornea has unmyelinated nerve endings sensitive to touch, temperature and chemicals; a touch of the cornea causes an involuntary reflex to close the eyelid. Because transparency is of prime importance, the healthy cornea does not have or need blood vessels within it. Instead, oxygen dissolves in tears and then diffuses throughout the cornea to keep it healthy. Similarly, nutrients are transported via diffusion from the tear fluid through the outside surface and the aqueous humour through the inside surface. Nutrients also come via neurotrophins supplied by the nerves of the cornea. In humans, the cornea has a diameter of about 11.5 mm and a thickness of 0.5–0.6 mm in the center and 0.6–0.8 mm at the periphery. Transparency, avascularity, the presence of immature resident immune cells, and immunologic privilege makes the cornea a very special tissue. The most abundant soluble protein in mammalian cornea is albumin. The human cornea borders with the sclera at the corneal limbus. In lampreys, the cornea is solely an extension of the sclera, and is separate from the skin above it, but in more advanced vertebrates it is always fused with the skin to form a single structure, albeit one composed of multiple layers. In fish, and aquatic vertebrates in general, the cornea plays no role in focusing light, since it has virtually the same refractive index as water. Microanatomy The human cornea has five layers (possibly six, if the Dua's layer is included). Corneas of other primates have five known layers. The corneas of cats, dogs, wolves, and other carnivores only have four. From the anterior to posterior the layers of the human cornea are: Corneal epithelium: an exceedingly thin multicellular epithelial tissue layer (non-keratinized stratified squamous epithelium) of fast-growing and easily regenerated cells, kept moist with tears. Irregularity or edema of the corneal epithelium disrupts the smoothness of the air/tear-film interface, the most significant component of the total refractive power of the eye, thereby reducing visual acuity. Corneal epithelium is continuous with the conjunctival epithelium, and is composed of about 6 layers of cells which are shed constantly on the exposed layer and are regenerated by multiplication in the basal layer. Bowman's layer (also known as the anterior limiting membrane): when discussed in lieu of a subepithelial basement membrane, Bowman's Layer is a tough layer composed of collagen (mainly type I collagen fibrils), laminin, nidogen, perlecan and other HSPGs that protects the corneal stroma. When discussed as a separate entity from the subepithelial basement membrane, Bowman's Layer can be described as an acellular, condensed region of the apical stroma, composed primarily of randomly organized yet tightly woven collagen fibrils. These fibrils interact with and attach onto each other. This layer is eight to 14 micrometres (μm) thick and is absent or very thin in non-primates. Corneal stroma (also substantia propria): a thick, transparent middle layer, consisting of regularly arranged collagen fibers along with sparsely distributed interconnected keratocytes, which are the cells for general repair and maintenance. They are parallel and are superimposed like book pages. The corneal stroma consists of approximately 200 layers of mainly type I collagen fibrils. Each layer is 1.5-2.5 μm. Up to 90% of the corneal thickness is composed of stroma. There are 2 theories of how transparency in the cornea comes about: The lattice arrangements of the collagen fibrils in the stroma. The light scatter by individual fibrils is cancelled by destructive interference from the scattered light from other individual fibrils. The spacing of the neighboring collagen fibrils in the stroma must be < 200 nm for there to be transparency. (Goldman and Benedek) Descemet's membrane (also posterior limiting membrane): a thin acellular layer that serves as the modified basement membrane of the corneal endothelium, from which the cells are derived. This layer is composed mainly of collagen type IV fibrils, less rigid than collagen type I fibrils, and is around 5-20 μm thick, depending on the subject's age. Just anterior to Descemet's membrane, a very thin and strong layer, Dua's layer, 15 microns thick and able to withstand 1.5 to 2 bars of pressure. Corneal endothelium: a simple squamous or low cuboidal monolayer, approx 5 μm thick, of mitochondria-rich cells. These cells are responsible for regulating fluid and solute transport between the aqueous and corneal stromal compartments. (The term endothelium is a misnomer here. The corneal endothelium is bathed by aqueous humor, not by blood or lymph, and has a very different origin, function, and appearance from vascular endothelia.) Unlike the corneal epithelium, the cells of the endothelium do not regenerate. Instead, they stretch to compensate for dead cells which reduces the overall cell density of the endothelium, which affects fluid regulation. If the endothelium can no longer maintain a proper fluid balance, stromal swelling due to excess fluids and subsequent loss of transparency will occur and this may cause corneal edema and interference with the transparency of the cornea and thus impairing the image formed. Iris pigment cells deposited on the corneal endothelium can sometimes be washed into a distinct vertical pattern by the aqueous currents - this is known as Krukenberg's Spindle. Nerve supply The cornea is one of the most sensitive tissues of the body, as it is densely innervated with sensory nerve fibres via the ophthalmic division of the trigeminal nerve by way of 70–80 long ciliary nerves. Research suggests the density of pain receptors in the cornea is 300–600 times greater than skin and 20–40 times greater than dental pulp, making any injury to the structure excruciatingly painful. The ciliary nerves run under the endothelium and exit the eye through holes in the sclera apart from the optic nerve (which transmits only optic signals). The nerves enter the cornea via three levels; scleral, episcleral and conjunctival. Most of the bundles give rise by subdivision to a network in the stroma, from which fibres supply the different regions. The three networks are, midstromal, subepithelial/sub-basal, and epithelial. The receptive fields of each nerve ending are very large, and may overlap. Corneal nerves of the subepithelial layer terminate near the superficial epithelial layer of the cornea in a logarithmic spiral pattern. The density of epithelial nerves decreases with age, especially after the seventh decade. Function Refraction The optical component is concerned with producing a reduced inverted image on the retina. The eye's optical system consists of not only two but four surfaces—two on the cornea, two on the lens. Rays are refracted toward the midline. Distant rays, due to their parallel nature, converge to a point on the retina. The cornea admits light at the greatest angle. The aqueous and vitreous humors both have a refractive index of 1.336-1.339, whereas the cornea has a refractive index of 1.376. Because the change in refractive index between cornea and aqueous humor is relatively small compared to the change at the air–cornea interface, it has a negligible refractive effect, typically -6 dioptres. The cornea is considered to be a positive meniscus lens. Some species of birds and chameleons, and one kinown species of fish, also have corneas which can focus. Transparency Upon death or removal of an eye the cornea absorbs the aqueous humor, thickens, and becomes hazy. Transparency can be restored by putting it in a warm, well-ventilated chamber at 31 °C (88 °F, the normal temperature), allowing the fluid to leave the cornea and become transparent. The cornea takes in fluid from the aqueous humor and the small blood vessels of the limbus, but a pump ejects the fluid immediately upon entry. When energy is deficient the pump may fail, or function too slowly to compensate, leading to swelling. This arises at death, but a dead eye can be placed in a warm chamber with a reservoir of sugar and glycogen that generally keeps the cornea transparent for at least 24 hours. The endothelium controls this pumping action, and as discussed above, damage thereof is more serious, and is a cause of opaqueness and swelling. When damage to the cornea occurs, such as in a viral infection, the collagen used to repair the process is not regularly arranged, leading to an opaque patch (leukoma). Clinical significance The most common corneal disorders are the following: Corneal abrasion – a medical condition involving the loss of the surface epithelial layer of the eye's cornea as a result of trauma to the surface of the eye. Corneal dystrophy – a condition in which one or more parts of the cornea lose their normal clarity due to a buildup of cloudy material. Corneal ulcer – an inflammatory or infective condition of the cornea involving disruption of its epithelial layer with involvement of the corneal stroma. Corneal neovascularization – excessive ingrowth of blood vessels from the limbal vascular plexus into the cornea, caused by deprivation of oxygen from the air. Fuchs' dystrophy – cloudy morning vision. Keratitis – inflammation of the cornea. Keratoconus – a degenerative disease, the cornea thins and changes shape to be more like a cone. Corneal foreign body – a foreign object present in the cornea, one of the most common preventable occupational hazards. Management Surgical procedures Various refractive eye surgery techniques change the shape of the cornea in order to reduce the need for corrective lenses or otherwise improve the refractive state of the eye. In many of the techniques used today, reshaping of the cornea is performed by photoablation using the excimer laser. There are also synthetic corneas (keratoprostheses) in development. Most are merely plastic inserts, but there are also those composed of biocompatible synthetic materials that encourage tissue ingrowth into the synthetic cornea, thereby promoting biointegration. Other methods, such as magnetic deformable membranes and optically coherent transcranial magnetic stimulation of the human retina are still in very early stages of research. Other procedures Orthokeratology is a method using specialized hard or rigid gas-permeable contact lenses to transiently reshape the cornea in order to improve the refractive state of the eye or reduce the need for eyeglasses and contact lenses. In 2009, researchers at the University of Pittsburgh Medical center demonstrated that stem cell collected from human corneas can restore transparency without provoking a rejection response in mice with corneal damage. For corneal epithelial diseases such as Stevens Johnson Syndrome, persistent corneal ulcer etc., the autologous contralateral (normal) suprabasal limbus derived in vitro expanded corneal limbal stem cells are found to be effective as amniotic membrane based expansion is controversial. For endothelial diseases, such as bullous keratopathy, cadaver corneal endothelial precursor cells have been proven to be efficient. Recently emerging tissue engineering technologies are expected to be capable of making one cadaver-donor's corneal cells be expanded and be usable in more than one patient's eye. Corneal retention and permeability in topical drug delivery to the eye The majority of ocular therapeutic agents are administered to the eye via the topical route. Cornea is one of the main barriers for drug diffusion because of its highly impermeable nature. Its continuous irrigation with a tear fluid also results in poor retention of the therapeutic agents on the ocular surface. Poor permeability of the cornea and quick wash out of therapeutic agents from ocular surface result in very low bioavailability of the drugs administered via topical route (typically less than 5%). Poor retention of formulations on ocular surfaces could potentially be improved with the use of mucoadhesive polymers. Drug permeability through the cornea could be facilitated with addition of penetration enhancers into topical formulations. Transplantation If the corneal stroma develops visually significant opacity, irregularity, or edema, a cornea of a deceased donor can be transplanted. Because there are no blood vessels in the cornea, there are also few problems with rejection of the new cornea. When a cornea is needed for transplant, as from an eye bank, the best procedure is to remove the cornea from the eyeball, preventing the cornea from absorbing the aqueous humor. There is a global shortage of corneal donations, severely limiting the availability of corneal transplants across most of the world. A 2016 study found that 12.7 million visually impaired people were in need of a corneal transplant, with only 1 cornea available for every 70 needed. Many countries have years-long waitlists for corneal transplant surgery due to the shortage of donated corneas. Only a handful of countries consistently have a large enough supply of donated corneas to meet local demand without a waitlist, including the United States, Italy, and Sri Lanka.	Biology and health sciences	Visual system	Biology
312008	https://en.wikipedia.org/wiki/Fourier%20optics	Fourier optics	Fourier optics is the study of classical optics using Fourier transforms (FTs), in which the waveform being considered is regarded as made up of a combination, or superposition, of plane waves. It has some parallels to the Huygens–Fresnel principle, in which the wavefront is regarded as being made up of a combination of spherical wavefronts (also called phasefronts) whose sum is the wavefront being studied. A key difference is that Fourier optics considers the plane waves to be natural modes of the propagation medium, as opposed to Huygens–Fresnel, where the spherical waves originate in the physical medium. A curved phasefront may be synthesized from an infinite number of these "natural modes" i.e., from plane wave phasefronts oriented in different directions in space. When an expanding spherical wave is far from its sources, it is locally tangent to a planar phase front (a single plane wave out of the infinite spectrum), which is transverse to the radial direction of propagation. In this case, a Fraunhofer diffraction pattern is created, which emanates from a single spherical wave phase center. In the near field, no single well-defined spherical wave phase center exists, so the wavefront isn't locally tangent to a spherical ball. In this case, a Fresnel diffraction pattern would be created, which emanates from an extended source, consisting of a distribution of (physically identifiable) spherical wave sources in space. In the near field, a full spectrum of plane waves is necessary to represent the Fresnel near-field wave, even locally. A "wide" wave moving forward (like an expanding ocean wave coming toward the shore) can be regarded as an infinite number of "plane wave modes", all of which could (when they collide with something such as a rock in the way) scatter independently of one other. These mathematical simplifications and calculations are the realm of Fourier analysis and synthesis – together, they can describe what happens when light passes through various slits, lenses or mirrors that are curved one way or the other, or is fully or partially reflected. Fourier optics forms much of the theory behind image processing techniques, as well as applications where information needs to be extracted from optical sources such as in quantum optics. To put it in a slightly complex way, similar to the concept of frequency and time used in traditional Fourier transform theory, Fourier optics makes use of the spatial frequency domain (kx, ky) as the conjugate of the spatial (x, y) domain. Terms and concepts such as transform theory, spectrum, bandwidth, window functions and sampling from one-dimensional signal processing are commonly used. Fourier optics plays an important role for high-precision optical applications such as photolithography in which a pattern on a reticle to be imaged on wafers for semiconductor chip production is so dense such that light (e.g., DUV or EUV) emanated from the reticle is diffracted and each diffracted light may correspond to a different spatial frequency (kx, ky). Due to generally non-uniform patterns on reticles, a simple diffraction grating analysis may not provide the details of how light is diffracted from each reticle. Propagation of light in homogeneous, source-free media Light can be described as a waveform propagating through a free space (vacuum) or a material medium (such as air or glass). Mathematically, a real-valued component of a vector field describing a wave is represented by a scalar wave function u that depends on both space and time: where represents a position in a three dimensional space (in the Cartesian coordinate system here), and t represents time. The wave equation Fourier optics begins with the homogeneous, scalar wave equation (valid in source-free regions): where is the speed of light and u(r,t) is a real-valued Cartesian component of an electromagnetic wave propagating through a free space (e.g., for where Ei is the i-axis component of an electric field E in the Cartesian coordinate system). Sinusoidal steady state If light of a fixed frequency in time/wavelength/color (as from a single-mode laser) is assumed, then, based on the engineering time convention, which assumes an time dependence in wave solutions at the angular frequency with where is a time period of the waves, the time-harmonic form of the optical field is given as where is the imaginary unit, is the operator taking the real part of , is the angular frequency (in radians per unit time) of light waves, and is, in general, a complex quantity, with separate amplitude in non-negative real number and phase . The Helmholtz equation Substituting this expression into the scalar wave equation above yields the time-independent form of the wave equation, where with the wavelength in vacuum, is the wave number (also called propagation constant), is the spatial part of a complex-valued Cartesian component of an electromagnetic wave. Note that the propagation constant and the angular frequency are linearly related to one another, a typical characteristic of transverse electromagnetic (TEM) waves in homogeneous media. Since the originally desired real-valued solution of the scalar wave equation can be simply obtained by taking the real part of , solving the following equation, known as the Helmholtz equation, is mostly concerned as treating a complex-valued function is often much easier than treating the corresponding real-valued function. Solving the Helmholtz equation Solutions to the Helmholtz equation in the Cartesian coordinate system may readily be found via the principle of separation of variables for partial differential equations. This principle says that in separable orthogonal coordinates, an elementary product solution to this wave equation may be constructed of the following form: i.e., as the product of a function of x, times a function of y, times a function of z. If this elementary product solution is substituted into the wave equation, using the scalar Laplacian in the Cartesian coordinates system then the following equation for the 3 individual functions is obtained which is readily rearranged into the form: It may now be argued that each quotient in the equation above must, of necessity, be constant. To justify this, let's say that the first quotient is not a constant, and is a function of x. Since none of the other terms in the equation has any dependence on the variable x, so the first term also must not have any x-dependence; it must be a constant. (If the first term is a function of x, then there is no way to make the left hand side of this equation be zero.) This constant is denoted as -kx2. Reasoning in a similar way for the y and z quotients, three ordinary differential equations are obtained for the fx, fy and fz, along with one separation condition: Each of these 3 differential equations has the same solution form: sines, cosines or complex exponentials. We'll go with the complex exponential as to be a complex function. As a result, the elementary product solution is with a generally complex number . This solution is the spatial part of a complex-valued Cartesian component (e.g., , , or as the electric field component along each axis in the Cartesian coordinate system) of a propagating plane wave. (, , or ) is a real number here since waves in a source-free medium has been assumed so each plane wave is not decayed or amplified as it propagates in the medium. The negative sign of (, , or ) in a wave vector (where ) means that the wave propagation direction vector has a positive (, , or )-component, while the positive sign of means a negative (, , or )-component of that vector. Product solutions to the Helmholtz equation are also readily obtained in cylindrical and spherical coordinates, yielding cylindrical and spherical harmonics (with the remaining separable coordinate systems being used much less frequently). The complete solution: the superposition integral A general solution to the homogeneous electromagnetic wave equation at a fixed time frequency in the Cartesian coordinate system may be formed as a weighted superposition of all possible elementary plane wave solutions as with the constraints of , each as a real number, and where . In this superposition, is the weight factor or the amplitude of the plane wave component with the wave vector where is determined in terms of and by the mentioned constraint. Next, let Then: The plane wave spectrum representation of a general electromagnetic field (e.g., a spherical wave) in the equation () is the basic foundation of Fourier optics (this point cannot be emphasized strongly enough), because at z = 0, the equation simply becomes a Fourier transform (FT) relationship between the field and its plane wave contents (hence the name, Fourier optics). Thus: and All spatial dependence of each plane wave component is described explicitly by an exponential function. The coefficient of the exponential is a function of only two components of the wave vector for each plane wave (since other remained component can be determined via the above mentioned constraints), for example and , just as in ordinary Fourier analysis and Fourier transforms. Connection between Fourier optics and imaging resolution Let's consider an imaging system where the z-axis is the optical axis of the system and the object plane (to be imaged on the image plane of the system) is the plane at . On the object plane, the spatial part of a complex-valued Cartesian component of a wave is, as shown above, with the constraints of , each as a real number, and where . The imaging is the reconstruction of a wave on the object plane (having information about a pattern on the object plane to be imaged) on the image plane via the proper wave propagation from the object to the image planes, (E.g., think about the imaging of an image in an aerial space.) and the wave on the object plane, that fully follows the pattern to be imaged, is in principle, described by the unconstrained inverse Fourier transform where takes an infinite range of real numbers. It means that, for a given light frequency, only a part of the full feature of the pattern can be imaged because of the above-mentioned constraints on ; (1) a fine feature which representation in the inverse Fourier transform requires spatial frequencies , where are transverse wave numbers satisfying , can not be fully imaged since waves with such do not exist for the given light of (This phenomenon is known as the diffraction limit.), and (2) spatial frequencies with but close to so higher wave outgoing angles with respect to the optical axis, requires a high NA (Numerical Aperture) imaging system that is expensive and difficult to build. For (1), even if complex-valued longitudinal wavenumbers are allowed (by an unknown interaction between light and the object plane pattern that is usually a solid material), give rise to light decay along the axis (Light amplification along the axis does not physically make sense if there is no amplification material between the object and image planes, and this is a usual case.) so waves with such may not reach the image plane that is usually sufficiently far way from the object plane. In connection with photolithography of electronic components, these (1) and (2) are the reasons why light of a higher frequency (smaller wavelength, thus larger magnitude of ) or a higher NA imaging system is required to image finer features of integrated circuits on a photoresist on a wafer. As a result, machines realizing such an optical lithography have become more and more complex and expensive, significantly increasing the cost of the electronic component production. The paraxial approximation Paraxial wave propagation (optic axis assumed as z axis) A solution to the Helmholtz equation as the spatial part of a complex-valued Cartesian component of a single frequency wave is assumed to take the form: where is the wave vector, and and is the wave number. Next, use the paraxial approximation, that is a small-angle approximation such that so, up to the second order approximation of trigonometric functions (that is, taking only up to the second term in the Taylor series expansion of each trigonometric function), where is the angle (in radian) between the wave vector k and the z-axis as the optical axis of an optical system under discussion. As a result, and The paraxial wave equation Substituting this expression into the Helmholtz equation, the paraxial wave equation is derived: where is the transverse Laplace operator in the Cartesian coordinates system. In the derivation of the paraxial wave equation, the following approximations are used. is small () so a term with is ignored. Terms with and are much smaller than a term with (or ) so these two terms are ignored. so a term with is ignored. It is the slowly varying envelope approximation, means that the amplitude or envelope of a wave is slowly varying compared with the major period of the wave . The far field approximation The equation () above may be evaluated asymptotically in the far field (using the stationary phase method) to show that the field at a distant point is indeed due solely to the plane wave component with the wave vector which propagates parallel to the vector , and whose plane is tangent to the phasefront at . The mathematical details of this process may be found in Scott [1998] or Scott [1990]. The result of performing a stationary phase integration on the expression above is the following expression, which clearly indicates that the field at is directly proportional to the spectral component in the direction of , where, and Stated another way, the radiation pattern of any planar field distribution is the FT (Fourier Transform) of that source distribution (see Huygens–Fresnel principle, wherein the same equation is developed using a Green's function approach). Note that this is NOT a plane wave. The radial dependence is a spherical wave - both in magnitude and phase - whose local amplitude is the FT of the source plane distribution at that far field angle. A plane wave spectrum does not necessarily mean that the field as the superposition of the plane wave components in that spectrum behaves something like a plane wave at far distances. Spatial versus angular bandwidth The equation () above is critical to making the connection between spatial bandwidth (on the one hand) and angular bandwidth (on the other), in the far field. Note that the term "far field" usually means we're talking about a converging or diverging spherical wave with a pretty well defined phase center. The connection between spatial and angular bandwidth in the far field is essential in understanding the low pass filtering property of thin lenses. See the section 6.1.3 for the condition defining the far field region. Once the concept of angular bandwidth is understood, the optical scientist can "jump back and forth" between the spatial and spectral domains to quickly gain insights which would ordinarily not be so readily available just through spatial domain or ray optics considerations alone. For example, any source bandwidth which lies past the edge angle to the first lens (This edge angle sets the bandwidth of the optical system.) will not be captured by the system to be processed. As a side note, electromagnetics scientists have devised an alternative means to calculate an electric field in a far zone which does not involve stationary phase integration. They have devised a concept known as "fictitious magnetic currents" usually denoted by M, and defined as In this equation, it is assumed that the unit vector in the z-direction points into the half-space where the far field calculations will be made. These equivalent magnetic currents are obtained using equivalence principles which, in the case of an infinite planar interface, allow any electric currents J to be "imaged away" while the fictitious magnetic currents are obtained from twice the aperture electric field (see Scott [1998]). Then the radiated electric field is calculated from the magnetic currents using an equation similar to the equation for the magnetic field radiated by an electric current. In this way, a vector equation is obtained for the radiated electric field in terms of the aperture electric field, and the derivation requires no use of stationary phase ideas. The plane wave spectrum: the foundation of Fourier optics The plane wave spectrum concept is the basic foundation of Fourier Optics. The plane wave spectrum is a continuous spectrum of uniform plane waves, and there is one plane wave component in the spectrum for every tangent point on the far-field phase front. The amplitude of that plane wave component would be the amplitude of the optical field at that tangent point. Again, this is true only in the far field, roughly defined as the range beyond where is the maximum linear extent of the optical sources and is the wavelength (Scott [1998]). The plane wave spectrum is often regarded as being discrete for certain types of periodic gratings, though in reality, the spectra from gratings are continuous as well, since no physical device can have the infinite extent required to produce a true line spectrum. Likely to electrical signals, bandwidth in optics is a measure of how finely detailed an image is; the finer the detail, the greater the bandwidth required to represent it. A DC (Direct Current) electrical signal is constant and has no oscillations; a plane wave propagating parallel to the optic () axis has constant value in any x-y plane, and therefore is analogous to the (constant) DC component of an electrical signal. Bandwidth in electrical signals relates to the difference between the highest and lowest frequencies present in the spectrum of a signal, practically with a criterion to cut off high and low frequency edges of the spectrum to represent bandwidth in a number. For optical systems, bandwidth also relates to spatial frequency content (spatial bandwidth), but it also has a secondary meaning. It also measures how far from the optic axis the corresponding plane waves are tilted, and so this type of bandwidth is often referred to also as angular bandwidth. It takes more frequency bandwidth to produce a short pulse in an electrical circuit, and more angular (or, spatial frequency) bandwidth to produce a sharp spot in an optical system (see discussion related to Point spread function). The plane wave spectrum arises naturally as the eigenfunction or "natural mode" solution to the homogeneous electromagnetic wave equation in rectangular coordinates (see also Electromagnetic radiation, which derives the wave equation from Maxwell's equations in source-free media, or Scott [1998]). In the frequency domain, with an assumed time convention of , the homogeneous electromagnetic wave equation becomes what is known as the Helmholtz equation and takes the form where and is the wavenumber of the medium. Eigenfunction (natural mode) solutions: background and overview In the case of differential equations, as in the case of matrix equations, whenever the right-hand side of an equation is zero (For example, a forcing function, forcing vector, or the source of a force is zero.), the equation may still admit a non-trivial solution, known in applied mathematics as an eigenfunction solution, in physics as a "natural mode" solution, and in electrical circuit theory as the "zero-input response." This is a concept that spans a wide range of physical disciplines. Common physical examples of resonant natural modes would include the resonant vibrational modes of stringed instruments (1D), percussion instruments (2D) or the former Tacoma Narrows Bridge (3D). Examples of propagating natural modes would include waveguide modes, optical fiber modes, solitons and Bloch waves. an Infinite homogeneous media admits the rectangular, circular and spherical harmonic solutions to the Helmholtz equation, depending on the coordinate system under consideration. The propagating plane waves that we'll study in this article are perhaps the simplest type of propagating waves found in any type of media. There is a striking similarity between the Helmholtz equation () above, which may be written and the usual equation form for the eigenvalues / eigenvectors of a square matrix A, particularly since both the scalar Laplacian and the matrix A are linear operators on their respective functions / vector spaces. (The minus sign in this matrix equation is, for all intents and purposes, immaterial. However, the plus sign in the Helmholtz equation is significant.) It is perhaps worthwhile to note that the eigenfunction solutions / eigenvector solutions to the Helmholtz equation / the matrix equation, often yield an orthogonal set of the eigenfunctions / the eigenvectors which span (i.e., form a basis set for) the function space / vector space under consideration. The interested reader may investigate other functional linear operators (so for different equations than the Helmholtz equation) which give rise to different kinds of orthogonal eigenfunctions such as Legendre polynomials, Chebyshev polynomials and Hermite polynomials. In the matrix equation case in which A is a square matrix, eigenvalues may be found by setting the determinant of the matrix equal to zero, i.e. finding where the matrix has no inverse. (Such a square matrix is said to be singular.) Finite matrices have only a finite number of eigenvalues/eigenvectors, whereas linear operators can have a countably infinite number of eigenvalues/eigenfunctions (in confined regions) or uncountably infinite (continuous) spectra of solutions, as in unbounded regions. In certain physics applications such as in the computation of bands in a periodic volume, it is often a case that the elements of a matrix will be very complicated functions of frequency and wavenumber, and the matrix will be non-singular (I.e., it has the inverse matrix.) for most combinations of frequency and wavenumber, but will also be singular (I.e., it does not have the inverse matrix.) for certain specific combinations. By finding which combinations of frequency and wavenumber drive the determinant of the matrix to zero, the propagation characteristics of the medium may be determined. Relations of this type, between frequency and wavenumber, are known as dispersion relations and some physical systems may admit many different kinds of dispersion relations. An example from electromagnetics is an ordinary waveguide, which may admit numerous dispersion relations, each associated with a unique propagation mode of the waveguide. Each propagation mode of the waveguide is known as an eigenfunction solution (or eigenmode solution) to Maxwell's equations in the waveguide. Free space also admits eigenmode (natural mode) solutions (known more commonly as plane waves), but with the distinction that for any given frequency, free space admits a continuous modal spectrum, whereas waveguides have a discrete mode spectrum. In this case, the dispersion relation is linear, as in section 1.3. K-space For a given such as for a homogeneous vacuum space, the separation condition, which is identical to the equation for the Euclidean metric in a three-dimensional configuration space, suggests the notion of a k-vector in a three-dimensional "k-space", defined (for propagating plane waves) in rectangular coordinates as: and in the spherical coordinate system as Use will be made of these spherical coordinate system relations in the next section. The notion of k-space is central to many disciplines in engineering and physics, especially in the study of periodic volumes, such as in crystallography and the band theory of semiconductor materials. The two-dimensional Fourier transform A spectrum analysis equation (calculating the spectrum of a function ): A synthesis equation (reconstructing the function from its spectrum): The normalizing factor of is present whenever angular frequency (radians) is used, but not when ordinary frequency (cycles) is used. Optical systems: general overview and analogy with electrical signal processing systems In a high level overview, an optical system consists of three parts; an input plane, and output plane, and a set of components between these planes that transform an image f formed in the input plane into a different image g formed in the output plane. The optical system output image g is related to the input image f by convolving the input image with the optical impulse response function of the optical system, h (known as the point-spread function, for focused optical systems). The impulse response function uniquely defines the input-output behavior of the optical system. By convention, the optical axis of the system is taken as the z-axis. As a result, the two images and the impulse response function are all functions of the transverse coordinates, x and y. The impulse response of an optical imaging system is the output plane field which is produced when an ideal mathematical optical field point source of light, that is an impulse input to the system, is placed in the input plane (usually on-axis, i.e., on the optical axis). In practice, it is not necessary to have an ideal point source in order to determine an exact impulse response. This is because any source bandwidth which lies outside the bandwidth of the optical system under consideration won't matter anyway (since it cannot even be captured by the optical system), so therefore it's not necessary in determining the impulse response. The source only needs to have at least as much (angular) bandwidth as the optical system. Optical systems typically fall into one of two different categories. The first is ordinary focused optical imaging systems (e.g., cameras), wherein the input plane is called the object plane and the output plane is called the image plane. An optical field in the image plane (the output plane of the imaging system) is desired to be a high-quality reproduction of an optical field in the object plane (the input plane of the imaging system). The impulse response function of an optical imaging system is desired to approximate a 2D delta function, at the location (or a linearly scaled location) in the output plane corresponding to the location of the impulse (an ideal point source) in the input plane. The actual impulse response function of an imaging system typically resembles an Airy function, whose radius is on the order of the wavelength of the light used. The impulse response function in this case is typically referred to as a point spread function, since the mathematical point of light in the object plane has been spread out into an Airy function in the image plane. The second type is optical image processing systems, in which a significant feature in the input plane optical field is to be located and isolated. In this case, the impulse response of such a system is desired to be a close replica (picture) of that feature which is being searched for in the input plane field, so that a convolution of the impulse response (an image of the desired feature) against the input plane field will produce a bright spot at the feature location in the output plane. It is this latter type of optical image processing system that is the subject of this section. The section 6.2 presents one hardware implementation of the optical image processing operations described in this section. Input plane The input plane is defined as the locus of all points such that z = 0. The input image f is therefore Output plane The output plane is defined as the locus of all points such that z = d. The output image g is therefore The 2D convolution of input function against the impulse response function i.e., The alert reader will note that the integral above tacitly assumes that the impulse response is NOT a function of the position (x',y') of the impulse of light in the input plane (if this were not the case, this type of convolution would not be possible). This property is known as shift invariance (Scott [1998]). No optical system is perfectly shift invariant: as the ideal, mathematical point of light is scanned away from the optic axis, aberrations will eventually degrade the impulse response (known as a coma in focused imaging systems). However, high quality optical systems are often "shift invariant enough" over certain regions of the input plane that we may regard the impulse response as being a function of only the difference between input and output plane coordinates, and thereby use the equation above with impunity. Also, this equation assumes unit magnification. If magnification is present, then eqn. () becomes which basically translates the impulse response function, hM(), from x′ to x = Mx′. In eqn. (), hM will be a magnified version of the impulse response function h of a similar, unmagnified system, so that hM(x,y) = h(x/M,y/M). Derivation of the convolution equation The extension to two dimensions is trivial, except for the difference that causality exists in the time domain, but not in the spatial domain. Causality means that the impulse response h(t − t′) of an electrical system, due to an impulse applied at time t', must of necessity be zero for all times t such that t − t′ < 0. Obtaining the convolution representation of the system response requires representing the input signal as a weighted superposition over a train of impulse functions by using the sifting property of Dirac delta functions. It is then presumed that the system under consideration is linear, that is to say that the output of the system due to two different inputs (possibly at two different times) is the sum of the individual outputs of the system to the two inputs, when introduced individually. Thus the optical system may contain no nonlinear materials nor active devices (except possibly, extremely linear active devices). The output of the system, for a single delta function input is defined as the impulse response of the system, h(t − t′). And, by our linearity assumption (i.e., that the output of system to a pulse train input is the sum of the outputs due to each individual pulse), we can now say that the general input function f(t) produces the output: where h(t − t′) is the (impulse) response of the linear system to the delta function input δ(t − t′), applied at time t. This is where the convolution equation above comes from. The convolution equation is useful because it is often much easier to find the response of a system to a delta function input - and then perform the convolution above to find the response to an arbitrary input - than it is to try to find the response to the arbitrary input directly. Also, the impulse response (in either time or frequency domains) usually yields insight to relevant figures of merit of the system. In the case of most lenses, the point spread function (PSF) is a pretty common figure of merit for evaluation purposes. The same logic is used in connection with the Huygens–Fresnel principle, or Stratton-Chu formulation, wherein the "impulse response" is referred to as the Green's function of the system. So the spatial domain operation of a linear optical system is analogous in this way to the Huygens–Fresnel principle. System transfer function If the last equation above is Fourier transformed, it becomes: where is the spectrum of the output signal is the system transfer function is the spectrum of the input signal In like fashion, eqn. () may be Fourier transformed to yield: The system transfer function, . In optical imaging this function is better known as the optical transfer function (Goodman). Once again it may be noted from the discussion on the Abbe sine condition, that this equation assumes unit magnification. This equation takes on its real meaning when the Fourier transform, is associated with the coefficient of the plane wave whose transverse wavenumbers are . Thus, the input-plane plane wave spectrum is transformed into the output-plane plane wave spectrum through the multiplicative action of the system transfer function. It is at this stage of understanding that the previous background on the plane wave spectrum becomes invaluable to the conceptualization of Fourier optical systems. Applications of Fourier optics principles Fourier optics is used in the field of optical information processing, the staple of which is the classical 4F processor. The Fourier transform properties of a lens provide numerous applications in optical signal processing such as spatial filtering, optical correlation and computer generated holograms. Fourier optical theory is used in interferometry, optical tweezers, atom traps, and quantum computing. Concepts of Fourier optics are used to reconstruct the phase of light intensity in the spatial frequency plane (see adaptive-additive algorithm). Fourier transforming property of lenses If a transmissive object is placed at one focal length in front of a lens, then its Fourier transform will be formed at one focal length behind the lens. Consider the figure to the right (click to enlarge) In this figure, a plane wave incident from the left is assumed. The transmittance function in the front focal plane (i.e., Plane 1) spatially modulates the incident plane wave in magnitude and phase, like on the left-hand side of eqn. () (specified to z = 0), and in so doing, produces a spectrum of plane waves corresponding to the FT of the transmittance function, like on the right-hand side of eqn. () (for z > 0). The various plane wave components propagate at different tilt angles with respect to the optic axis of the lens (i.e., the horizontal axis). The finer the features in the transparency, the broader the angular bandwidth of the plane wave spectrum. We'll consider one such plane wave component, propagating at angle θ with respect to the optic axis. It is assumed that θ is small (paraxial approximation), so that and and In the figure, the plane wave phase, moving horizontally from the front focal plane to the lens plane, is and the spherical wave phase from the lens to the spot in the back focal plane is: and the sum of the two path lengths is f (1 + θ2/2 + 1 − θ2/2) = 2f; i.e., it is a constant value, independent of tilt angle, θ, for paraxial plane waves. Each paraxial plane wave component of the field in the front focal plane appears as a point spread function spot in the back focal plane, with an intensity and phase equal to the intensity and phase of the original plane wave component in the front focal plane. In other words, the field in the back focal plane is the Fourier transform of the field in the front focal plane. All FT components are computed simultaneously - in parallel - at the speed of light. As an example, light travels at a speed of roughly per nanosecond, so if a lens has a focal length, an entire 2D FT can be computed in about 2 ns (2 × 10−9 seconds). If the focal length is 1 in, then the time is under 200 ps. No electronic computer can compete with these kinds of numbers or perhaps ever hope to, although supercomputers may actually prove faster than optics, as improbable as that may seem. However, their speed is obtained by combining numerous computers which, individually, are still slower than optics. The disadvantage of the optical FT is that, as the derivation shows, the FT relationship only holds for paraxial plane waves, so this FT "computer" is inherently bandlimited. On the other hand, since the wavelength of visible light is so minute in relation to even the smallest visible feature dimensions in the image i.e., (for all kx, ky within the spatial bandwidth of the image, so that kz is nearly equal to k), the paraxial approximation is not terribly limiting in practice. And, of course, this is an analog - not a digital - computer, so precision is limited. Also, phase can be challenging to extract; often it is inferred interferometrically. Optical processing is especially useful in real time applications where rapid processing of massive amounts of 2D data is required, particularly in relation to pattern recognition. Object truncation and Gibbs phenomenon The spatially modulated electric field, shown on the left-hand side of eqn. (), typically only occupies a finite (usually rectangular) aperture in the x,y plane. The rectangular aperture function acts like a 2D square-top filter, where the field is assumed to be zero outside this 2D rectangle. The spatial domain integrals for calculating the FT coefficients on the right-hand side of eqn. () are truncated at the boundary of this aperture. This step truncation can introduce inaccuracies in both theoretical calculations and measured values of the plane wave coefficients on the RHS of eqn. (). Whenever a function is discontinuously truncated in one FT domain, broadening and rippling are introduced in the other FT domain. A perfect example from optics is in connection with the point spread function, which for on-axis plane wave illumination of a quadratic lens (with circular aperture), is an Airy function, J1(x)/x. Literally, the point source has been "spread out" (with ripples added), to form the Airy point spread function (as the result of truncation of the plane wave spectrum by the finite aperture of the lens). This source of error is known as Gibbs phenomenon and it may be mitigated by simply ensuring that all significant content lies near the center of the transparency, or through the use of window functions which smoothly taper the field to zero at the frame boundaries. By the convolution theorem, the FT of an arbitrary transparency function - multiplied (or truncated) by an aperture function - is equal to the FT of the non-truncated transparency function convolved against the FT of the aperture function, which in this case becomes a type of "Greens function" or "impulse response function" in the spectral domain. Therefore, the image of a circular lens is equal to the object plane function convolved against the Airy function (the FT of a circular aperture function is J1(x)/x and the FT of a rectangular aperture function is a product of sinc functions, sinx/x). Fourier analysis and functional decomposition Even though the input transparency only occupies a finite portion of the x-y plane (Plane 1), the uniform plane waves comprising the plane wave spectrum occupy the entire x-y plane, which is why (for this purpose) only the longitudinal plane wave phase (in the z-direction, from Plane 1 to Plane 2) must be considered, and not the phase transverse to the z-direction. It is of course, very tempting to think that if a plane wave emanating from the finite aperture of the transparency is tilted too far from horizontal, it will somehow "miss" the lens altogether but again, since the uniform plane wave extends infinitely far in all directions in the transverse (x-y) plane, the planar wave components cannot miss the lens. This issue brings up perhaps the predominant difficulty with Fourier analysis, namely that the input-plane function, defined over a finite support (i.e., over its own finite aperture), is being approximated with other functions (sinusoids) which have infinite support (i.e., they are defined over the entire infinite x-y plane). This is unbelievably inefficient computationally, and is the principal reason why wavelets were conceived, that is to represent a function (defined on a finite interval or area) in terms of oscillatory functions which are also defined over finite intervals or areas. Thus, instead of getting the frequency content of the entire image all at once (along with the frequency content of the entire rest of the x-y plane, over which the image has zero value), the result is instead the frequency content of different parts of the image, which is usually much simpler. Unfortunately, wavelets in the x-y plane don't correspond to any known type of propagating wave function, in the same way that Fourier's sinusoids (in the x-y plane) correspond to plane wave functions in three dimensions. However, the FTs of most wavelets are well known and could possibly be shown to be equivalent to some useful type of propagating field. On the other hand, sinc functions and Airy functions - which are not only the point spread functions of rectangular and circular apertures, respectively, but are also cardinal functions commonly used for functional decomposition in interpolation/sampling theory [Scott 1990] - do''' correspond to converging or diverging spherical waves, and therefore could potentially be implemented as a whole new functional decomposition of the object plane function, thereby leading to another point of view similar in nature to Fourier optics. This would basically be the same as conventional ray optics, but with diffraction effects included. In this case, each point spread function would be a type of "smooth pixel," in much the same way that a soliton on a fiber is a "smooth pulse." Perhaps a lens figure-of-merit in this "point spread function" viewpoint would be to ask how well a lens transforms an Airy function in the object plane into an Airy function in the image plane, as a function of radial distance from the optic axis, or as a function of the size of the object plane Airy function. This is somewhat like the point spread function, except now we're really looking at it as a kind of input-to-output plane transfer function (like MTF), and not so much in absolute terms, relative to a perfect point. Similarly, Gaussian wavelets, which would correspond to the waist of a propagating Gaussian beam, could also potentially be used in still another functional decomposition of the object plane field. Far-field range and the 2D2 / λ criterion In the figure above, illustrating the Fourier transforming property of lenses, the lens is in the near field of the object plane transparency, therefore the object plane field at the lens may be regarded as a superposition of plane waves, each one of which propagates at some angle with respect to the z-axis. In this regard, the far-field criterion is loosely defined as: Range = 2D2/λ where D is the maximum linear extent of the optical sources and λ is the wavelength (Scott [1998]). The D of the transparency is on the order of cm (10−2 m) and the wavelength of light is on the order of 10−6 m, therefore D/λ for the whole transparency is on the order of 104. This times D is on the order of 102 m, or hundreds of meters. On the other hand, the far field distance from a PSF spot is on the order of λ. This is because D for the spot is on the order of λ, so that D/λ is on the order of unity; this times D (i.e., λ) is on the order of λ (10−6 m). Since the lens is in the far field of any PSF spot, the field incident on the lens from the spot may be regarded as being a spherical wave, as in eqn. (), not as a plane wave spectrum, as in eqn. (). On the other hand, the lens is in the near field of the entire input plane transparency, therefore eqn. () - the full plane wave spectrum - accurately represents the field incident on the lens from that larger, extended source. Lens as a low-pass filter A lens is basically a low-pass plane wave filter (see Low-pass filter). Consider a "small" light source located on-axis in the object plane of the lens. It is assumed that the source is small enough that, by the far-field criterion, the lens is in the far field of the "small" source. Then, the field radiated by the small source is a spherical wave which is modulated by the FT of the source distribution, as in eqn. (), Then, the lens passes - from the object plane over onto the image plane - only that portion of the radiated spherical wave which lies inside the edge angle of the lens. In this far-field case, truncation of the radiated spherical wave is equivalent to truncation of the plane wave spectrum of the small source. So, the plane wave components in this far-field spherical wave, which lie beyond the edge angle of the lens, are not captured by the lens and are not transferred over to the image plane. Note: this logic is valid only for small sources, such that the lens is in the far field region of the source, according to the 2D2/λ criterion mentioned previously. If an object plane transparency is imagined as a summation over small sources (as in the Whittaker–Shannon interpolation formula, Scott [1990]), each of which has its spectrum truncated in this fashion, then every point of the entire object plane transparency suffers the same effects of this low pass filtering. Loss of the high (spatial) frequency content causes blurring and loss of sharpness (see discussion related to point spread function). Bandwidth truncation causes a (fictitious, mathematical, ideal) point source in the object plane to be blurred (or, spread out) in the image plane, giving rise to the term, "point spread function." Whenever bandwidth is expanded or contracted, image size is typically contracted or expanded accordingly, in such a way that the space-bandwidth product remains constant, by Heisenberg's principle (Scott [1998] and Abbe sine condition). Coherence and Fourier transforming While working in the frequency domain, with an assumed ejωt (engineering) time dependence, coherent (laser) light is implicitly assumed, which has a delta function dependence in the frequency domain. Light at different (delta function) frequencies will "spray" the plane wave spectrum out at different angles, and as a result these plane wave components will be focused at different places in the output plane. The Fourier transforming property of lenses works best with coherent light, unless there is some special reason to combine light of different frequencies, to achieve some special purpose. Hardware implementation of the system transfer function: The 4F correlator The theory on optical transfer functions presented in the section 5 is somewhat abstract. However, there is one very well known device which implements the system transfer function H in hardware using only 2 identical lenses and a transparency plate - the 4F correlator. Although one important application of this device would certainly be to implement the mathematical operations of cross-correlation and convolution, this device - 4 focal lengths long - actually serves a wide variety of image processing operations that go well beyond what its name implies. A diagram of a typical 4F correlator is shown in the figure below (click to enlarge). This device may be readily understood by combining the plane wave spectrum representation of the electric field (section 1.5) with the Fourier transforming property of quadratic lenses (section 6.1) to yield the optical image processing operations described in the section 5. The 4F correlator is based on the convolution theorem from Fourier transform theory, which states that convolution in the spatial (x,y) domain is equivalent to direct multiplication in the spatial frequency (kx, ky) domain (aka: spectral domain). Once again, a plane wave is assumed incident from the left and a transparency containing one 2D function, f(x,y), is placed in the input plane of the correlator, located one focal length in front of the first lens. The transparency spatially modulates the incident plane wave in magnitude and phase, like on the left-hand side of eqn. (), and in so doing, produces a spectrum of plane waves corresponding to the FT of the transmittance function, like on the right-hand side of eqn. (). That spectrum is then formed as an "image" one focal length behind the first lens, as shown. A transmission mask containing the FT of the second function, g(x,y), is placed in this same plane, one focal length behind the first lens, causing the transmission through the mask to be equal to the product, F(kx,ky) × G(kx,ky). This product now lies in the "input plane" of the second lens (one focal length in front), so that the FT of this product (i.e., the convolution of f(x,y) and g(x,y)), is formed in the back focal plane of the second lens. If an ideal, mathematical point source of light is placed on-axis in the input plane of the first lens, then there will be a uniform, collimated field produced in the output plane of the first lens. When this uniform, collimated field is multiplied by the FT plane mask, and then Fourier transformed by the second lens, the output plane field (which in this case is the impulse response of the correlator) is just our correlating function, g(x,y). In practical applications, g(x,y) will be some type of feature which must be identified and located within the input plane field (see Scott [1998]). In military applications, this feature may be a tank, ship or airplane which must be quickly identified within some more complex scene. The 4F correlator is an excellent device for illustrating the "systems" aspects of optical instruments, alluded to in the section 5 above. The FT plane mask function, G(kx,ky) is the system transfer function of the correlator, which we'd in general denote as H(kx,ky), and it is the FT of the impulse response function of the correlator, h(x,y) which is just our correlating function g(x,y). And, as mentioned above, the impulse response of the correlator is just a picture of the feature we're trying to find in the input image. In the 4F correlator, the system transfer function H(kx,ky) is directly multiplied against the spectrum F(kx,ky) of the input function, to produce the spectrum of the output function. This is how electrical signal processing systems operate on 1D temporal signals. Image restoration Image blurring by a point spread function is studied extensively in optical information processing, one way to alleviate the blurring is to adopt Wiener Filter. For example, assume that is the intensity distribution from an incoherent object, is the intensity distribution of its image which is blurred by a space-invariant point-spread function and a noise introduced in the detection process: The goal of image restoration is to find a linear restoration filter that minimize the mean-squared error between the true distribution and the estimation . That is, to minimize The solution of this optimization problem is Wiener filter: where , , are the power spectral densities of the point-spread function, the object and the noise. Ragnarsson proposed a method to realize Wiener restoration filters optically by holographic technique like setup shown in the figure. The derivation of the function of the setup is described as follows. Assume there is a transparency as the recording plane and an impulse emitted from a point source S. The wave of impulse is collimated by lens L1, forming a distribution equal to the impulse response . Then the distribution is then split into two parts: The upper portion is first focused (i.e., Fourier transformed) by a lens L2 to a spot in the front focal plan of lens L3, forming a virtual point source generating a spherical wave. The wave is then collimated by lens L3 and produces a tilted plane wave with the form at the recording plane. The lower portion is directly collimated by lens L3'', yielding an amplitude distribution . Therefore, the total intensity distribution is Assume has an amplitude distribution and a phase distribution such that then we can rewrite intensity as follows: Note that for the point at the origin of the film plane (), the recorded wave from the lower portion should be much stronger than that from the upper portion because the wave passing through the lower path is focused, which leads to the relationship . In Ragnarsson' s work, this method is based on the following postulates: Assume there is a transparency, with its amplitude transmittance proportional to , that has recorded the known impulse response of the blurred system. The maximum phase shift introduced by the filter is much smaller than radians so that . The phase shift of the transparency after bleaching is linearly proportional to the silver density present before bleaching. The density is linearly proportional to the logarithm of exposure The average exposure is much stronger than varying exposure By these postulates, we have the following relationship: Finally, we get a amplitude transmittance with the form of a Wiener filter: Afterword: Plane wave spectrum within the broader context of functional decomposition Electrical fields can be represented mathematically in many different ways. In the Huygens–Fresnel or Stratton-Chu viewpoints, the electric field is represented as a superposition of point sources, each one of which gives rise to a Green's function field. The total field is then the weighted sum of all of the individual Green's function fields. That seems to be the most natural way of viewing the electric field for most people - no doubt because most of us have, at one time or another, drawn out the circles with protractor and paper, much the same way Thomas Young did in his classic paper on the double-slit experiment. However, it is by no means the only way to represent the electric field, which may also be represented as a spectrum of sinusoidally varying plane waves. In addition, Frits Zernike proposed still another functional decomposition based on his Zernike polynomials, defined on the unit disc. The third-order (and lower) Zernike polynomials correspond to the normal lens aberrations. And still another functional decomposition could be made in terms of Sinc functions and Airy functions, as in the Whittaker–Shannon interpolation formula and the Nyquist–Shannon sampling theorem. All of these functional decompositions have utility in different circumstances. The optical scientist having access to these various representational forms has available a richer insight to the nature of these marvelous fields and their properties. These different ways of looking at the field are not conflicting or contradictory, rather, by exploring their connections, one can often gain deeper insight into the nature of wave fields. Functional decomposition and eigenfunctions The twin subjects of eigenfunction expansions and functional decomposition, both briefly alluded to here, are not completely independent. The eigenfunction expansions to certain linear operators defined over a given domain, will often yield a countably infinite set of orthogonal functions which will span that domain. Depending on the operator and the dimensionality (and shape, and boundary conditions) of its domain, many different types of functional decompositions are, in principle, possible.	Physical sciences	Optics	Physics
312018	https://en.wikipedia.org/wiki/VMware	VMware	VMware LLC is an American cloud computing and virtualization technology company headquartered in Palo Alto, California. VMware was the first commercially successful company to virtualize the x86 architecture. VMware's desktop software runs on Microsoft Windows, Linux, and macOS. VMware ESXi, its enterprise software hypervisor, is an operating system that runs on server hardware. On November 22, 2023, Broadcom Inc. acquired VMware in a cash-and-stock transaction valued at US$69 billion, with the End-User Computing (EUC) division of VMware then sold to KKR and rebranded to Omnissa. History Early history In 1998, VMware was founded by Diane Greene, Mendel Rosenblum, Scott Devine, Ellen Wang, and Edouard Bugnion. Greene and Rosenblum were graduate students at the University of California, Berkeley. Edouard Bugnion remained the chief architect and CTO of VMware until 2005 and went on to found Nuova Systems (now part of Cisco). VMware operated in stealth mode for the first year, with roughly 20 employees by the end of 1998. The company was launched officially early in the second year, in February 1999, at the DEMO conference organized by Chris Shipley. The first product, VMware Workstation, was delivered in May 1999, and the company entered the server market in 2001 with VMware GSX Server (hosted) and VMware ESX Server (host-less). In 2003, VMware launched VMware Virtual Center, vMotion, and Virtual Symmetric Multi-Processing (SMP) technology. 64-bit support was introduced in 2004. Acquisition by EMC On January 9, 2004, under the terms of the definitive agreement announced on December 15, 2003, EMC (now Dell EMC) acquired the company with US$625 million in cash. On August 14, 2007, EMC sold 15% of VMware to the public via an initial public offering. Shares were priced at per share and closed the day at . On July 8, 2008, after disappointing financial performance, the board of directors fired VMware co-founder, president and CEO Diane Greene, who was replaced by Paul Maritz, a 14-year Microsoft veteran who was heading EMC's cloud computing business unit. Greene had been CEO since the company's founding, ten years earlier. On September 10, 2008, Mendel Rosenblum, the company's co-founder, chief scientist, and the husband of Diane Greene, resigned. On September 16, 2008, VMware announced a collaboration with Cisco Systems. One result was the Cisco Nexus 1000V, a distributed virtual software switch, an integrated option in the VMware infrastructure. In April 2011, EMC transferred control of the Mozy backup service to VMware. On April 12, 2011, VMware released an open-source platform-as-a-service system called Cloud Foundry, as well as a hosted version of the service. This supported application deployment for Java, Ruby on Rails, Sinatra, Node.js, and Scala, as well as database support for MySQL, MongoDB, Redis, PostgreSQL, and RabbitMQ. In August 2012, Pat Gelsinger was appointed as the new CEO of VMware, coming over from EMC. Paul Maritz went over to EMC as Head of Strategy before moving on to lead the Pivotal spin-off. In March 2013, VMware announced the corporate spin-off of Pivotal Software, with General Electric investing in the company. Most of VMware's application- and developer-oriented products, including Spring, tc Server, Cloud Foundry, RabbitMQ, GemFire, and SQLFire were transferred to this organization. In May 2013, VMware launched its own IaaS service, vCloud Hybrid Service, at its new Palo Alto headquarters (vCloud Hybrid Service was rebranded vCloud Air and later sold to cloud provider OVH), announcing an early access program in a Las Vegas data center. The service is designed to function as an extension of its customer's existing vSphere installations, with full compatibility with existing virtual machines virtualized with VMware software and tightly integrated networking. The service is based on vCloud Director 5.1/vSphere 5.1. In September 2013, at VMworld San Francisco, VMware announced the general availability of vCloud Hybrid Service and expansion to Sterling, Virginia, Santa Clara, California, Dallas, Texas, and a service beta in the UK. It announced the acquisition of Desktone in October 2013. Acquisition by Dell In January 2016, in anticipation of Dell's acquisition of EMC, VMware announced a restructuring to reduce about 800 positions, and some executives resigned. The entire development team behind VMware Workstation and Fusion was disbanded and all US developers were immediately fired. On April 24, 2016, maintenance release 12.1.1 was released. On September 8, 2016, VMware announced the release of Workstation 12.5 and Fusion 8.5 as a free upgrade supporting Windows 10 and Windows Server 2016. In April 2016, VMware president and COO Carl Eschenbach left VMware to join Sequoia Capital, and Martin Casado, VMware's general manager for its Networking and Security business, left to join Andreessen Horowitz. Analysts commented that the cultures at Dell and EMC, and at EMC and VMware, are different, and said that they had heard that impending corporate cultural collisions and potentially radical product overlap pruning, would cause many EMC and VMware personnel to leave; VMware CEO Pat Gelsinger, following rumors, categorically denied that he would leave. In August 2016 VMware introduced the VMware Cloud Provider website. Mozy was transferred to Dell in 2016 after the merger of Dell and EMC. In April 2017, according to Glassdoor, VMware was ranked 3rd on the list of highest paying companies in the United States. In Q2 2017, VMware sold vCloud Air to French cloud service provider OVH. On January 13, 2021, VMware announced that CEO Pat Gelsinger would be leaving to step in at Intel. Intel is where Gelsinger spent 30 years of his career and was Intel's first chief technology officer. CFO Zane Rowe became interim CEO while the board searched for a replacement. Spinoff from Dell On April 15, 2021, it was reported that Dell would spin off its remaining stake in VMware to shareholders and that the two companies would continue to operate without major changes for at least five years. The spinoff was completed on November 1, 2021. On May 12, 2021, VMware announced that Raghu Raghuram would take over as CEO. In May 2022, VMware announced that the company had partnered with Formula One motor racing team, McLaren Racing. Log4Shell vulnerability Beginning in January 2022, hackers infiltrated servers using the Log4Shell vulnerability at organizations who failed to implement available patches released by VMware according to PCMag. ZDNET reported in March 2022 that hackers utilized Log4Shell on some customers' VMware servers to install backdoors and for cryptocurrency mining. In May 2022, Bleeping Computer reported that the Lazarus Group cybercrime group, which is possibly linked to North Korea, was actively using Log4Shell "to inject backdoors that fetch information-stealing payloads on VMware Horizon servers", including VMware Horizon. Acquisition by Broadcom On May 26, 2022, Broadcom announced its intention to acquire VMware for approximately $61 billion in cash and stock in addition to assuming $8 billion of VMware's net debt, and that Broadcom Software Group would rebrand and operate as VMware. In November 2022, the UK's Competition and Markets Authority regulator announced it would investigate whether the acquisition would "result in a substantial lessening of competition within any market or markets in the United Kingdom for goods or services". The transaction closed on November 22, 2023, after a prolonged delay in getting approval from the Chinese regulator on an additional condition that VMware's server software should maintain compatibility with third-party hardware and not require the use of Broadcom's hardware products. On completion, Broadcom reorganized the company into four divisions: VMware Cloud Foundation, Tanzu, Software-Defined Edge, and Application Networking and Security, and subsequently laid off over 2,800 employees. Broadcom also relocated its headquarters from North San Jose to VMware's headquarters campus in Palo Alto. On December 13, 2023, VMware ended availability for perpetually licensed products such as vSphere and Cloud Foundation, moving exclusively to subscription-based offerings. The company stated that this had been planned as an eventuality prior to the Broadcom acquisition. In February 2024 private equity firm KKR and Broadcom agreed for KKR to acquire Broadcom's End-User Computing (EUC) Division, formerly a division of VMware, for about $4 billion. This includes the VDI product Horizon and the device management suite Workspace ONE, formerly AirWatch. On May 14, 2024, it was announced that VMware Workstation Pro and VMware Fusion Pro would be made free for personal use, with commercial use still requiring payment. In November 2024, VMware announced that commercial use would be free too. Acquisitions Litigation In March 2015, the Software Freedom Conservancy announced it was funding litigation by Christoph Hellwig in Hamburg, Germany against VMware for alleged violation of his copyrights in its ESXi product. Hellwig's core claim is that ESXi is a derivative work of the GPLv2-licensed Linux kernel 2.4, and therefore VMware is not in compliance with GPLv2 because it does not publish the source code to ESXi. VMware publicly stated that ESXi is not a derivative of the Linux kernel, denying Hellwig's core claim. VMware said it offered a way to use Linux device drivers with ESXi, and that code does use some Linux GPLv2-licensed code and so it had published the source, meeting GPLv2 requirements. The lawsuit was dismissed by the court in July 2016 and Hellwig announced he would file an appeal. The appeal was decided February 2019 and again dismissed by German court, on the basis of not meeting "procedural requirements for the burden of proof of the plaintiff." In May 2023, VMware was ordered to pay $84.5 million for patent infringement on two patents belonging to Densify, a Canadian software company. Current products VMware's most notable products are its hypervisors. VMware became well known for its first type 2 hypervisor known as VMware Workstation. This product has since evolved into two additional hypervisor product lines: VMware's type 1 hypervisors running directly on hardware (ESX/ESXi) and their discontinued hosted type 2 hypervisors (GSX). VMware software provides a completely virtualized set of hardware to the guest operating system. VMware software virtualizes the hardware for a video adapter, a network adapter, and hard disk adapters. The host provides pass-through drivers for guest USB, serial, and parallel devices. In this way, VMware virtual machines become highly portable between computers, because every host looks nearly identical to the guest. In practice, a system administrator can pause operations on a virtual machine guest, move or copy that guest to another physical computer, and there resume execution exactly at the point of suspension. Alternatively, for enterprise servers, a feature called vMotion allows the migration of operational guest virtual machines between similar but separate hardware hosts sharing the same storage (or, with vMotion Storage, separate storage can be used, too). Each of these transitions is completely transparent to any users on the virtual machine at the time it is being migrated. VMware's products predate the virtualization extensions to the x86 instruction set, and do not require virtualization-enabled processors. On newer processors, the hypervisor is now designed to take advantage of the extensions. However, unlike many other hypervisors, VMware still supports older processors. In such cases, it uses the CPU to run code directly whenever possible (as, for example, when running user-mode and virtual 8086 mode code on x86). When direct execution cannot operate, such as with kernel-level and real-mode code, VMware products use binary translation (BT) to re-write the code dynamically. The translated code gets stored in spare memory, typically at the end of the address space, which segmentation mechanisms can protect and make invisible. For these reasons, VMware operates dramatically faster than emulators, running at more than 80% of the speed that the virtual guest operating system would run directly on the same hardware. In one study VMware claims a slowdown over native ranging from 0–6 percent for the VMware ESX Server. Desktop software VMware Workstation, introduced in 1999, was the first product launched by VMware. This software suite allows users to run multiple instances of x86 or x86-64-compatible operating systems on a single physical personal computer. Version 17.0 was released on November 17, 2022. Originally a commercial app, WMware Workstation has become freeware in December 2024. VMware Fusion (discontinued on 30 April 2024), provides similar functionality for users of the Intel Mac platform, along with full compatibility with virtual machines created by other VMware products. VMware Workstation Player (discontinued) was freeware for non-commercial use, without requiring a license, and available for commercial use with permission. It is similar to VMware Workstation, with some features not available, including support for UEFI Secure Boot, snapshots, encrypted virtual machines, and some advanced features. Server software VMware ESXi, an enterprise software product, can deliver greater performance than the freeware VMware Server, due to lower system computational overhead. VMware ESXi, as a "bare-metal" product, runs directly on the server hardware, allowing virtual servers to also use hardware more or less directly. In addition, VMware ESXi integrates into VMware vCenter, which offers extra services. Cloud management software VMware Suite – a cloud management platform purpose-built for a hybrid cloud. VMware vRealize Hyperic was acquired from SpringSource and subsequently discontinued in 2020. VMware Go is a web-based service to guide users of any expertise level through the installation and configuration of VMware vSphere Hypervisor. VMware Cloud Foundation – Cloud Foundation provides an easy way to deploy and operate a private cloud on an integrated SDDC system. VMware Horizon View is a virtual desktop infrastructure (VDI) product. vSphere+ and vSAN+ – activates add-on hybrid cloud services for business-critical applications running on-premises, including IT disaster recovery and ransomware protection Application management VMware Workspace Portal was a self-service app store for workspace management. Provisioning PlateSpin (does Provisioning) Storage and availability VMware's storage and availability products are composed of two primary offerings: VMware vSAN (previously called VMware Virtual SAN) is software-defined storage that is embedded in VMware's ESXi hypervisor. The vSphere and vSAN software runs on industry-standard x86 servers to form a hyper-converged infrastructure (or HCI). However, network operators need to have servers from HCL (Hardware Compatibility List) to put one into production. The first release, version 5.5, was released in March 2014. The 6th generation, version 6.6, was released in April 2017. New features available in VMware vSAN 6.6 include native data at rest encryption, local protection for stretched clusters, analytics, and optimized solid-state drive performance. The VMWare 6.7 version was released in April 2018. VMware Site Recovery Manager (SRM) automates the failover and failback of virtual machines to and from a secondary site using policy-based management. Networking and security products VMware NSX is VMware's network virtualization product marketed using the term software-defined data center (SDDC). The technology included some acquired from the 2012 purchase of Nicira. Software Defined Networking (SDN) allows the same policies that govern Identity and Access Management (IAM) to dictate levels of access to applications and data through a totally converged infrastructure not possible with legacy network and system access methods. Other products Workspace ONE allows mobile users to access apps and data. The VIX (Virtual Infrastructure eXtension) API allows automated or scripted management of a computer virtualized using either VMware's vSphere, Workstation, Player, or Fusion products. VIX provides bindings for the programming languages C, Perl, Visual Basic, VBScript and C#. Herald is a communications protocol from VMware for more reliable Bluetooth communication and range finding for mobile devices. Herald code is available under an open-source license and was implemented in the Australian Government's COVIDSafe app for contact tracing on 19 December 2020.	Technology	Virtualization	null
312029	https://en.wikipedia.org/wiki/Cowrie	Cowrie	Cowrie or cowry () is the common name for a group of small to large sea snails in the family Cypraeidae. The term porcelain derives from the old Italian term for the cowrie shell () due to their similar appearance. Cowrie shells have held cultural, economic, and ornamental significance in various cultures. The cowrie was the shell most widely used worldwide as shell money. It is most abundant in the Indian Ocean, and was collected in the Maldive Islands, in Sri Lanka, along the Indian Malabar coast, in Borneo and on other East Indian islands, in Maluku in the Pacific, and in various parts of the African coast from Ras Hafun to Mozambique. Cowrie shell money was important in the trade networks of Africa, South Asia, and East Asia. In the United States and Mexico, cowrie species inhabit the waters off Central California to Baja California (the chestnut cowrie is the only cowrie species native to the eastern Pacific Ocean off the coast of the United States; further south, off the coast of Mexico, Central America and Peru, Little Deer Cowrie habitat can be found; and further into the Pacific from Central America, the Pacific habitat range of Money Cowrie can be reached) as well as the waters south of the Southeastern United States. Some species in the family Ovulidae are also often referred to as cowries. In the British Isles the local Trivia species (family Triviidae, species Trivia monacha and Trivia arctica) are sometimes called cowries. The Ovulidae and the Triviidae are other families within Cypraeoidea, the superfamily of cowries and their close relatives. Etymology The word cowrie comes from Hindi (), which is itself derived from Sanskrit (). Shell description The shells of cowries are usually smooth and shiny and more or less egg-shaped. The round side of the shell is called the Dorsal Face, whereas the flat under side is called the Ventral Face, which shows a long, narrow, slit-like opening (aperture), which is often toothed at the edges. The narrower end of the egg-shaped cowrie shell is the anterior end, and the broader end of the shell is called the posterior. The spire of the shell is not visible in the adult shell of most species, but is visible in juveniles, which have a different shape from the adults. Nearly all cowries have a porcelain-like shine, with some exceptions such as Hawaii's granulated cowrie, Nucleolaria granulata. Many have colorful patterns. Lengths range from for some species up to for the Atlantic deer cowrie, Macrocypraea cervus. Human use Monetary use Cowrie shells, especially Monetaria moneta, were used for centuries as currency by native Africans. In his book Marriage and Morals, Bertrand Russell attributed the use of cowrie shells as currency in ancient Egypt to the similarity between shape of the shell and that of female genitalia. Additionally, the money cowrie was almost impossible to counterfeit until the late 19th Century. After the 1500s, however, the shell's use as currency became even more common. Western nations, chiefly through the slave trade, introduced huge numbers of Maldivian cowries in Africa. The Ghanaian cedi was named after cowrie shells. Starting over three thousand years ago, cowrie shells, or copies of the shells, were used as Chinese currency. They were also used as means of exchange in India. The Classical Chinese character for money (貝) originated as a stylized drawing of a Maldivian cowrie shell. Words and characters concerning money, property or wealth usually have this as a radical. Before the Spring and Autumn period the cowrie was used as a type of trade token awarding access to a feudal lord's resources to a worthy vassal. Ritual use The Ojibwe aboriginal people in North America use cowrie shells which are called sacred miigis shells or whiteshells in Midewiwin ceremonies, and the Whiteshell Provincial Park in Manitoba, Canada is named after this type of shell. There is some debate about how the Ojibway traded for or found these shells, so far inland and so far north, very distant from the natural habitat. Oral stories and birch bark scrolls seem to indicate that the shells were found in the ground, or washed up on the shores of lakes or rivers. Finding the cowrie shells so far inland could indicate the previous use of them by an earlier tribe or group in the area, who may have obtained them through an extensive trade network in the ancient past. In Eastern India, particularly in West Bengal, it is given as a token price for the ferry ride of the departed soul to cross the river "Vaitarani". Cowries are used during cremation. Cowries are also used in the worship of Goddess Laxmi. In Brazil, as a result of the Atlantic slave trade from Africa, cowrie shells (called búzios) are also used to consult the Orixás divinities and hear their replies. Cowrie shells were among the devices used for divination by the Kaniyar Panicker astrologers of Kerala, India. In certain parts of Africa, cowries were prized charms, and they were said to be associated with fecundity, sexual pleasure and good luck. It is also used in the treatment of certain diseases such as rashes and ringworm when it is burnt into ashes. In Pre-dynastic Egypt and Neolithic Southern Levant, cowrie shells were placed in the graves of young girls. The modified Levantine cowries were discovered ritually arranged around the skull in female burials. During the Bronze Age, cowries became more common as funerary goods, also associated with burials of women and children. Jewelry Cowrie shells are also worn as jewelry or otherwise used as ornaments or charms. In Mende culture, cowrie shells are viewed as symbols of womanhood, fertility, birth and wealth. Its underside is supposed, by one modern ethnographic author, to represent a vulva or an eye. On the Fiji Islands, a shell of the golden cowrie or bulikula, Cypraea aurantium, was drilled at the ends and worn on a string around the neck by chieftains as a badge of rank. The women of Tuvalu use cowrie and other shells in traditional handicrafts. Games and gambling Cowrie shells are sometimes used in a way similar to dice, e.g., in board games like Pachisi, Ashta Chamma or in divination (cf. Ifá and the annual customs of Dahomey of Benin). A number of shells (6 or 7 in Pachisi) are thrown, with those landing aperture upwards indicating the actual number rolled. In Nepal cowries are used for a gambling game, where 16 pieces of cowries are tossed by four different bettors (and sub-bettors under them). This game is usually played at homes and in public during the Hindu festival of Tihar or Deepawali. In the same festival these shells are also worshiped as a symbol of Goddess Lakshmi and wealth. Other Large cowrie shells such as that of a Cypraea tigris have been used in Europe in the recent past as a darning egg over which sock heels were stretched. The cowrie's smooth surface allows the needle to be positioned under the cloth more easily. In the 1940s and 1950s, small cowry shells were used as a teaching aid in infant schools e.g counting, adding, subtracting.	Biology and health sciences	Gastropods	Animals
312249	https://en.wikipedia.org/wiki/Bryophyte	Bryophyte	Bryophytes () are a group of land plants (embryophytes), sometimes treated as a taxonomic division, that contains three groups of non-vascular land plants: the liverworts, hornworts, and mosses. In the strict sense, the division Bryophyta consists of the mosses only. Bryophytes are characteristically limited in size and prefer moist habitats although some species can survive in drier environments. The bryophytes consist of about 20,000 plant species. Bryophytes produce enclosed reproductive structures (gametangia and sporangia), but they do not produce flowers or seeds. They reproduce sexually by spores and asexually by fragmentation or the production of gemmae. Though bryophytes were considered a paraphyletic group in recent years, almost all of the most recent phylogenetic evidence supports the monophyly of this group, as originally classified by Wilhelm Schimper in 1879. The term bryophyte comes . Features The defining features of bryophytes are: Their life cycles are dominated by a multicellular haploid gametophyte stage Their sporophytes are diploid and unbranched They do not have a true vascular tissue containing lignin (although some have specialized tissues for the transport of water) Habitat Bryophytes exist in a wide variety of habitats. They can be found growing in a range of temperatures (cold arctics and in hot deserts), elevations (sea-level to alpine), and moisture (dry deserts to wet rain forests). Bryophytes can grow where vascularized plants cannot because they do not depend on roots for uptake of nutrients from soil. Bryophytes can survive on rocks and bare soil. Life cycle Like all land plants (embryophytes), bryophytes have life cycles with alternation of generations. In each cycle, a haploid gametophyte, each of whose cells contains a fixed number of unpaired chromosomes, alternates with a diploid sporophyte, whose cells contain two sets of paired chromosomes. Gametophytes produce haploid sperm and eggs which fuse to form diploid zygotes that grow into sporophytes. Sporophytes produce haploid spores by meiosis, that grow into gametophytes. Bryophytes are gametophyte dominant, meaning that the more prominent, longer-lived plant is the haploid gametophyte. The diploid sporophytes appear only occasionally and remain attached to and nutritionally dependent on the gametophyte. In bryophytes, the sporophytes are always unbranched and produce a single sporangium (spore producing capsule), but each gametophyte can give rise to several sporophytes at once. Liverworts, mosses and hornworts spend most of their lives as gametophytes. Gametangia (gamete-producing organs), archegonia and antheridia, are produced on the gametophytes, sometimes at the tips of shoots, in the axils of leaves or hidden under thalli. Some bryophytes, such as the liverwort Marchantia, create elaborate structures to bear the gametangia that are called gametangiophores. Sperm are flagellated and must swim from the antheridia that produce them to archegonia which may be on a different plant. Arthropods can assist in transfer of sperm. Fertilized eggs become zygotes, which develop into sporophyte embryos inside the archegonia. Mature sporophytes remain attached to the gametophyte. They consist of a stalk called a seta and a single sporangium or capsule. Inside the sporangium, haploid spores are produced by meiosis. These are dispersed, most commonly by wind, and if they land in a suitable environment can develop into a new gametophyte. Thus bryophytes disperse by a combination of swimming sperm and spores, in a manner similar to lycophytes, ferns and other cryptogams. The sporophyte develops differently in the three groups. Both mosses and hornworts have a meristem zone where cell division occurs. In hornworts, the meristem starts at the base where the foot ends, and the division of cells pushes the sporophyte body upwards. In mosses, the meristem is located between the capsule and the top of the stalk (seta), and produces cells downward, elongating the stalk and elevating the capsule. In liverworts the meristem is absent and the elongation of the sporophyte is caused almost exclusively by cell expansion. Sexuality The arrangement of antheridia and archegonia on an individual bryophyte plant is usually constant within a species, although in some species it may depend on environmental conditions. The main division is between species in which the antheridia and archegonia occur on the same plant and those in which they occur on different plants. The term monoicous may be used where antheridia and archegonia occur on the same gametophyte and the term dioicous where they occur on different gametophytes. In seed plants, "monoecious" is used where flowers with anthers (microsporangia) and flowers with ovules (megasporangia) occur on the same sporophyte and "dioecious" where they occur on different sporophytes. These terms occasionally may be used instead of "monoicous" and "dioicous" to describe bryophyte gametophytes. "Monoecious" and "monoicous" are both derived from the Greek for "one house", "dioecious" and "dioicous" from the Greek for two houses. The use of the "-oicy" terminology refers to the gametophyte sexuality of bryophytes as distinct from the sporophyte sexuality of seed plants. Monoicous plants are necessarily hermaphroditic, meaning that the same plant produces gametes of both sexes. The exact arrangement of the antheridia and archegonia in monoicous plants varies. They may be borne on different shoots (autoicous), on the same shoot but not together in a common structure (paroicous or paroecious), or together in a common "inflorescence" (synoicous or synoecious). Dioicous plants are unisexual, meaning that an individual plant has only one sex. All four patterns (autoicous, paroicous, synoicous and dioicous) occur in species of the moss genus Bryum. Classification and phylogeny Traditionally, all living land plants without vascular tissues were classified in a single taxonomic group, often a division (or phylum). The term "Bryophyta" was first suggested by Braun in 1864. As early as 1879, the term Bryophyta was used by German bryologist Wilhelm Schimper to describe a group containing all three bryophyte clades (though at the time, hornworts were considered part of the liverworts). G.M. Smith placed this group between Algae and Pteridophyta. Although a 2005 study supported this traditional monophyletic view, by 2010 a broad consensus had emerged among systematists that bryophytes as a whole are not a natural group (i.e., are paraphyletic). However, a 2014 study concluded that these previous phylogenies, which were based on nucleic acid sequences, were subject to composition biases, and that, furthermore, phylogenies based on amino acid sequences suggested that the bryophytes are monophyletic after all. Since then, partially thanks to a proliferation of genomic and transcriptomic datasets, almost all phylogenetics studies based on nuclear and chloroplastic sequences have concluded that the bryophytes form a monophyletic group. Nevertheless, phylogenies based on mitochondrial sequences fail to support the monophyletic view. The three bryophyte clades are the Marchantiophyta (liverworts), Bryophyta (mosses) and Anthocerotophyta (hornworts). However, it has been proposed that these clades are de-ranked to the classes Marchantiopsida, Bryopsida, and Anthocerotopsida, respectively. There is now strong evidence that the liverworts and mosses belong to a monophyletic clade, called Setaphyta. Monophyletic view The favored model, based on amino acids phylogenies, indicates bryophytes as a monophyletic group: Consistent with this view, compared to other living land plants, all three lineages lack vascular tissue containing lignin and branched sporophytes bearing multiple sporangia. The prominence of the gametophyte in the life cycle is also a shared feature of the three bryophyte lineages (extant vascular plants are all sporophyte dominant). However, if this phylogeny is correct, then the complex sporophyte of living vascular plants might have evolved independently of the simpler unbranched sporophyte present in bryophytes. Furthermore, this view implies that stomata evolved only once in plant evolution, before being subsequently lost in the liverworts. Paraphyletic view In this alternative view, the Setaphyta grouping is retained, but hornworts instead are sister to vascular plants. (Another paraphyletic view involves hornworts branching out first.) Traditional morphology Traditionally, when basing classifications on morphological characters, bryophytes have been distinguished by their lack of vascular structure. However, this distinction is problematic, firstly because some of the earliest-diverging (but now extinct) non-bryophytes, such as the horneophytes, did not have true vascular tissue, and secondly because many mosses have well-developed water-conducting vessels. A more useful distinction may lie in the structure of their sporophytes. In bryophytes, the sporophyte is a simple unbranched structure with a single spore-forming organ (sporangium), whereas in all other land plants, the polysporangiophytes, the sporophyte is branched and carries many sporangia. The contrast is shown in the cladogram below: Evolution There have probably been several different terrestrialization events, in which originally aquatic organisms colonized the land, just within the lineage of the Viridiplantae. Between 510 and 630 million years ago, however, land plants emerged within the green algae. Molecular phylogenetic studies conclude that bryophytes are the earliest diverging lineages of the extant land plants. They provide insights into the migration of plants from aquatic environments to land. A number of physical features link bryophytes to both land plants and aquatic plants. Similarities to algae and vascular plants Green algae, bryophytes and vascular plants all have chlorophyll a and b, and the chloroplast structures are similar. Like green algae and land plants, bryophytes also produce starch stored in the plastids and contain cellulose in their walls. Distinct adaptations observed in bryophytes have allowed plants to colonize Earth's terrestrial environments. To prevent desiccation of plant tissues in a terrestrial environment, a waxy cuticle covering the soft tissue of the plant may be present, providing protection. In hornworts and mosses, stomata provide gas exchange between the atmosphere and an internal intercellular space system. The development of gametangia provided further protection specifically for gametes, the zygote and the developing sporophyte. The bryophytes and vascular plants (embryophytes) also have embryonic development which is not seen in green algae. While bryophytes have no truly vascularized tissue, they do have organs that are specialized for transport of water and other specific functions, analogous for example to the functions of leaves and stems in vascular land plants. Bryophytes depend on water for reproduction and survival. In common with ferns and lycophytes, a thin layer of water is required on the surface of the plant to enable the movement of the flagellated sperm between gametophytes and the fertilization of an egg. Comparative morphology Summary of the morphological characteristics of the gametophytes of the three groups of bryophytes: Summary of the morphological characteristics of the sporophytes of the three groups of bryophytes: Uses Environmental Characteristics of bryophytes make them useful to the environment. Depending on the specific plant texture, bryophytes have been shown to help improve the water retention and air space within soil. Bryophytes are used in pollution studies to indicate soil pollution (such as the presence of heavy metals), air pollution, and UV-B radiation. Gardens in Japan are designed with moss to create peaceful sanctuaries. Some bryophytes have been found to produce natural pesticides. The liverwort, Plagiochila, produces a chemical that is poisonous to mice. Other bryophytes produce chemicals that are antifeedants which protect them from being eaten by slugs. When Phythium sphagnum is sprinkled on the soil of germinating seeds, it inhibits growth of "damping off fungus" which would otherwise kill young seedlings. Commercial Peat is a fuel produced from dried bryophytes, typically Sphagnum. Bryophytes' antibiotic properties and ability to retain water make them a useful packaging material for vegetables, flowers, and bulbs. Also, because of its antibiotic properties, Sphagnum was used as a surgical dressing in World War I.	Biology and health sciences	Bryophytes	null
312252	https://en.wikipedia.org/wiki/Partition%20function%20%28statistical%20mechanics%29	Partition function (statistical mechanics)	In physics, a partition function describes the statistical properties of a system in thermodynamic equilibrium. Partition functions are functions of the thermodynamic state variables, such as the temperature and volume. Most of the aggregate thermodynamic variables of the system, such as the total energy, free energy, entropy, and pressure, can be expressed in terms of the partition function or its derivatives. The partition function is dimensionless. Each partition function is constructed to represent a particular statistical ensemble (which, in turn, corresponds to a particular free energy). The most common statistical ensembles have named partition functions. The canonical partition function applies to a canonical ensemble, in which the system is allowed to exchange heat with the environment at fixed temperature, volume, and number of particles. The grand canonical partition function applies to a grand canonical ensemble, in which the system can exchange both heat and particles with the environment, at fixed temperature, volume, and chemical potential. Other types of partition functions can be defined for different circumstances; see partition function (mathematics) for generalizations. The partition function has many physical meanings, as discussed in Meaning and significance. Canonical partition function Definition Initially, let us assume that a thermodynamically large system is in thermal contact with the environment, with a temperature T, and both the volume of the system and the number of constituent particles are fixed. A collection of this kind of system comprises an ensemble called a canonical ensemble. The appropriate mathematical expression for the canonical partition function depends on the degrees of freedom of the system, whether the context is classical mechanics or quantum mechanics, and whether the spectrum of states is discrete or continuous. Classical discrete system For a canonical ensemble that is classical and discrete, the canonical partition function is defined as where is the index for the microstates of the system; is Euler's number; is the thermodynamic beta, defined as where is the Boltzmann constant; is the total energy of the system in the respective microstate. The exponential factor is otherwise known as the Boltzmann factor. Classical continuous system In classical mechanics, the position and momentum variables of a particle can vary continuously, so the set of microstates is actually uncountable. In classical statistical mechanics, it is rather inaccurate to express the partition function as a sum of discrete terms. In this case we must describe the partition function using an integral rather than a sum. For a canonical ensemble that is classical and continuous, the canonical partition function is defined as where is the Planck constant; is the thermodynamic beta, defined as ; is the Hamiltonian of the system; is the canonical position; is the canonical momentum. To make it into a dimensionless quantity, we must divide it by h, which is some quantity with units of action (usually taken to be the Planck constant). Classical continuous system (multiple identical particles) For a gas of identical classical non-interacting particles in three dimensions, the partition function is where is the Planck constant; is the thermodynamic beta, defined as ; is the index for the particles of the system; is the Hamiltonian of a respective particle; is the canonical position of the respective particle; is the canonical momentum of the respective particle; is shorthand notation to indicate that and are vectors in three-dimensional space. is the classical continuous partition function of a single particle as given in the previous section. The reason for the factorial factor N! is discussed below. The extra constant factor introduced in the denominator was introduced because, unlike the discrete form, the continuous form shown above is not dimensionless. As stated in the previous section, to make it into a dimensionless quantity, we must divide it by h3N (where h is usually taken to be the Planck constant). Quantum mechanical discrete system For a canonical ensemble that is quantum mechanical and discrete, the canonical partition function is defined as the trace of the Boltzmann factor: where: is the trace of a matrix; is the thermodynamic beta, defined as ; is the Hamiltonian operator. The dimension of is the number of energy eigenstates of the system. Quantum mechanical continuous system For a canonical ensemble that is quantum mechanical and continuous, the canonical partition function is defined as where: is the Planck constant; is the thermodynamic beta, defined as ; is the Hamiltonian operator; is the canonical position; is the canonical momentum. In systems with multiple quantum states s sharing the same energy Es, it is said that the energy levels of the system are degenerate. In the case of degenerate energy levels, we can write the partition function in terms of the contribution from energy levels (indexed by j) as follows: where gj is the degeneracy factor, or number of quantum states s that have the same energy level defined by Ej = Es. The above treatment applies to quantum statistical mechanics, where a physical system inside a finite-sized box will typically have a discrete set of energy eigenstates, which we can use as the states s above. In quantum mechanics, the partition function can be more formally written as a trace over the state space (which is independent of the choice of basis): where is the quantum Hamiltonian operator. The exponential of an operator can be defined using the exponential power series. The classical form of Z is recovered when the trace is expressed in terms of coherent states and when quantum-mechanical uncertainties in the position and momentum of a particle are regarded as negligible. Formally, using bra–ket notation, one inserts under the trace for each degree of freedom the identity: where is a normalised Gaussian wavepacket centered at position x and momentum p. Thus A coherent state is an approximate eigenstate of both operators and , hence also of the Hamiltonian , with errors of the size of the uncertainties. If and can be regarded as zero, the action of reduces to multiplication by the classical Hamiltonian, and reduces to the classical configuration integral. Connection to probability theory For simplicity, we will use the discrete form of the partition function in this section. Our results will apply equally well to the continuous form. Consider a system S embedded into a heat bath B. Let the total energy of both systems be E. Let pi denote the probability that the system S is in a particular microstate, i, with energy Ei. According to the fundamental postulate of statistical mechanics (which states that all attainable microstates of a system are equally probable), the probability pi will be inversely proportional to the number of microstates of the total closed system (S, B) in which S is in microstate i with energy Ei. Equivalently, pi will be proportional to the number of microstates of the heat bath B with energy : Assuming that the heat bath's internal energy is much larger than the energy of S (), we can Taylor-expand to first order in Ei and use the thermodynamic relation , where here , are the entropy and temperature of the bath respectively: Thus Since the total probability to find the system in some microstate (the sum of all pi) must be equal to 1, we know that the constant of proportionality must be the normalization constant, and so, we can define the partition function to be this constant: Calculating the thermodynamic total energy In order to demonstrate the usefulness of the partition function, let us calculate the thermodynamic value of the total energy. This is simply the expected value, or ensemble average for the energy, which is the sum of the microstate energies weighted by their probabilities: or, equivalently, Incidentally, one should note that if the microstate energies depend on a parameter λ in the manner then the expected value of A is This provides us with a method for calculating the expected values of many microscopic quantities. We add the quantity artificially to the microstate energies (or, in the language of quantum mechanics, to the Hamiltonian), calculate the new partition function and expected value, and then set λ to zero in the final expression. This is analogous to the source field method used in the path integral formulation of quantum field theory. Relation to thermodynamic variables In this section, we will state the relationships between the partition function and the various thermodynamic parameters of the system. These results can be derived using the method of the previous section and the various thermodynamic relations. As we have already seen, the thermodynamic energy is The variance in the energy (or "energy fluctuation") is The heat capacity is In general, consider the extensive variable X and intensive variable Y where X and Y form a pair of conjugate variables. In ensembles where Y is fixed (and X is allowed to fluctuate), then the average value of X will be: The sign will depend on the specific definitions of the variables X and Y. An example would be X = volume and Y = pressure. Additionally, the variance in X will be In the special case of entropy, entropy is given by where A is the Helmholtz free energy defined as , where is the total energy and S is the entropy, so that Furthermore, the heat capacity can be expressed as Partition functions of subsystems Suppose a system is subdivided into N sub-systems with negligible interaction energy, that is, we can assume the particles are essentially non-interacting. If the partition functions of the sub-systems are ζ1, ζ2, ..., ζN, then the partition function of the entire system is the product of the individual partition functions: If the sub-systems have the same physical properties, then their partition functions are equal, ζ1 = ζ2 = ... = ζ, in which case However, there is a well-known exception to this rule. If the sub-systems are actually identical particles, in the quantum mechanical sense that they are impossible to distinguish even in principle, the total partition function must be divided by a N! (N factorial): This is to ensure that we do not "over-count" the number of microstates. While this may seem like a strange requirement, it is actually necessary to preserve the existence of a thermodynamic limit for such systems. This is known as the Gibbs paradox. Meaning and significance It may not be obvious why the partition function, as we have defined it above, is an important quantity. First, consider what goes into it. The partition function is a function of the temperature T and the microstate energies E1, E2, E3, etc. The microstate energies are determined by other thermodynamic variables, such as the number of particles and the volume, as well as microscopic quantities like the mass of the constituent particles. This dependence on microscopic variables is the central point of statistical mechanics. With a model of the microscopic constituents of a system, one can calculate the microstate energies, and thus the partition function, which will then allow us to calculate all the other thermodynamic properties of the system. The partition function can be related to thermodynamic properties because it has a very important statistical meaning. The probability Ps that the system occupies microstate s is Thus, as shown above, the partition function plays the role of a normalizing constant (note that it does not depend on s), ensuring that the probabilities sum up to one: This is the reason for calling Z the "partition function": it encodes how the probabilities are partitioned among the different microstates, based on their individual energies. Other partition functions for different ensembles divide up the probabilities based on other macrostate variables. As an example: the partition function for the isothermal-isobaric ensemble, the generalized Boltzmann distribution, divides up probabilities based on particle number, pressure, and temperature. The energy is replaced by the characteristic potential of that ensemble, the Gibbs Free Energy. The letter Z stands for the German word Zustandssumme, "sum over states". The usefulness of the partition function stems from the fact that the macroscopic thermodynamic quantities of a system can be related to its microscopic details through the derivatives of its partition function. Finding the partition function is also equivalent to performing a Laplace transform of the density of states function from the energy domain to the β domain, and the inverse Laplace transform of the partition function reclaims the state density function of energies. Grand canonical partition function We can define a grand canonical partition function for a grand canonical ensemble, which describes the statistics of a constant-volume system that can exchange both heat and particles with a reservoir. The reservoir has a constant temperature T, and a chemical potential μ. The grand canonical partition function, denoted by , is the following sum over microstates Here, each microstate is labelled by , and has total particle number and total energy . This partition function is closely related to the grand potential, , by the relation This can be contrasted to the canonical partition function above, which is related instead to the Helmholtz free energy. It is important to note that the number of microstates in the grand canonical ensemble may be much larger than in the canonical ensemble, since here we consider not only variations in energy but also in particle number. Again, the utility of the grand canonical partition function is that it is related to the probability that the system is in state : An important application of the grand canonical ensemble is in deriving exactly the statistics of a non-interacting many-body quantum gas (Fermi–Dirac statistics for fermions, Bose–Einstein statistics for bosons), however it is much more generally applicable than that. The grand canonical ensemble may also be used to describe classical systems, or even interacting quantum gases. The grand partition function is sometimes written (equivalently) in terms of alternate variables as where is known as the absolute activity (or fugacity) and is the canonical partition function.	Physical sciences	Statistical mechanics	Physics
312259	https://en.wikipedia.org/wiki/Chinampa	Chinampa	Chinampa ( ) is a technique used in Mesoamerican agriculture which relies on small, rectangular areas of fertile arable land to grow crops on the shallow lake beds in the Valley of Mexico. The word chinampa has Nahuatl origins, chinampa meaning “in the fence of reeds”. They are built up on wetlands of a lake or freshwater swamp for agricultural purposes, and their proportions ensure optimal moisture retention. This method was also used and occupied most of Lake Xochimilco. The United Nations designated it a Globally Important Agricultural Heritage System in 2018. Although different technologies existed during the Post-classic and Colonial periods in the basin, chinampas have raised many questions on agricultural production and political development. After the Aztec Triple Alliance formed, the conquest of southern basin city-states, such as Xochimilco, was one of the first strategies of imperial expansion. Before this time, farmers maintained small-scale chinampas adjacent to their households and communities in the freshwater lakes of Xochimilco and Chalco. The Aztecs did not invent chinampas but rather were the first to develop it to a large scale cultivation. Sometimes referred to as "floating gardens," chinampas are artificial islands that were created by interweaving reeds with stakes beneath the lake's surface, creating underwater fences. A buildup of soil and aquatic vegetation would be piled into these "fences" until the top layer of soil was visible on the water's surface. When creating chinampas, in addition to building up masses of land, a drainage system was developed. This drainage system was multi-purposed. A ditch was created to allow for the flow of water and sediments (likely including night soil). Over time, the ditch would slowly accumulate piles of mud. This mud would then be dug up and placed on top of the chinampas, clearing the blockage. The soil from the bottom of the lake was also rich in nutrients, thus acting as an efficient and effective way of fertilizing the chinampas. Replenishing the topsoil with lost nutrients provided for bountiful harvests. Embarcadero-Jiménez and colleagues tested the correlation between environmental parameters and bacterial diversity in the soil. It is speculated that a diverse array of bacteria can affect the nutrients in the soil. The results found that bacterial diversity was more abundant in cultivated soils than non-cultivated soils. Also, "the structure of the bacterial communities showed that the chinampas are a transition system between sediment and soil and revealed an interesting association of the S-cycle and iron-oxidizing bacteria with the rhizosphere of plants grown in the chinampa soil". Evidence from Nahuatl wills from late seventeenth-century Pueblo Culhuacán suggests chinampas were measured in matl (one matl = 1.67 meters), often listed in groups of seven. One scholar has calculated the size of chinampas using Codex Vergara as a source, finding that they usually measured roughly . In Tenochtitlan, the chinampas ranged from to They were created by staking out the shallow lake bed and then fencing in the rectangle with wattle. The fenced-off area was then layered with mud, lake sediment, and decaying vegetation, eventually bringing it above the level of the lake. Often trees such as āhuexōtl (Salix bonplandiana) (a willow) and āhuēhuētl (Taxodium mucronatum) (a cypress) were planted at the corners to secure the chinampa. In some places, the long raised beds had ditches in between them, giving plants continuous access to water and making crops grown there independent of rainfall. Chinampas were separated by channels wide enough for a canoe to pass. These raised, well-watered beds had very high crop yields with up to 7 harvests a year. Chinampas were commonly used in pre-colonial Mexico and Central America. There is evidence that the Nahua settlement of Culhuacan, on the south side of the Ixtapalapa peninsula that divided Lake Texcoco from Lake Xochimilco, constructed the first chinampas in C.E. 1100. History The earliest fields that have been securely dated are from the Middle Postclassic period, 1150 – 1350 CE. Chinampas were used primarily in Lakes Xochimilco and Chalco near the springs that lined the south shore of those lakes. The Aztecs not only conducted military campaigns to obtain control over these regions but, according to some researchers, undertook significant state-led efforts to increase their extent. There is some strong evidence to suggest state-led operations for the “expansion” of the chinampas. This is sometimes referred to as the hydraulic hypothesis, which is directly related to a hydraulic empire, which is an empire that maintains power and control through the regulation and distribution of water. There is evidence to support the idea of state involvement, primarily the amount of manpower and materials it would take to build, turn, and maintain the chinampas. However, arguments about state control of the chinampas rely upon the assumption that dikes were necessary to control the water levels and to keep the saline water of Lake Texcoco away from the freshwater of the chinampa zone. This is plausible, but there is evidence that the chinampas were functional before the construction of a dike that protected them from the saline water. It is suggested that the dike was meant to drastically improve the size of the chinampa operation. Chinampa farms also ringed Tenochtitlán, the Aztec capital, which was considerably enlarged over time. Smaller-scale farms have also been identified near the island-city of Xaltocan and on the east side of Lake Texcoco. With the destruction of the dams and sluice gates during the Spanish conquest of the Aztec Empire, many chinampas fields were abandoned. However, many lakeshore towns retained their chinampas through the end of the colonial era since cultivation was highly labor-intensive and less attractive for Spaniards to acquire. The Aztecs built Tenochtitlan on an island around 1325. Issues arose when the cities' constant expansion eventually caused them to run out of room to build. As the empire grew, more sources of food were required. At times this meant conquering more land; at other times it meant expanding the chinampa system. With this expansion, chinampas' multiple crops per year became a large factor in the production and supply of food. Empirical records suggest that farmers had a relatively light tribute to pay compared to others because the annual tribute may have been only a fraction of the amount necessary for local needs. The extent to which Tenochtitlan depended on chinampas for its fresh food supply has been the topic of a number of scholarly studies. Among the crops grown on chinampas were maize, beans, squash, amaranth, tomatoes, chili peppers, and flowers. Maize was planted with digging stick huictli with a wooden blade on one end. The word chinampa comes from the Nahuatl word chināmitl, meaning "square made of canes" and the Nahuatl locative, "pan." In documentation by Spaniards, they used the word camellones, "ridges between the rows." However, Franciscan Fray Juan de Torquemada described them with the Nahua term, chinampa, saying "without much trouble [the Indians] plant and harvest their maize and greens, for all over there are ridges called chinampas; these were strips built above water and surrounded by ditches, which obviates watering." Chinampas are depicted in pictorial Aztec codices, including Codex Vergara, Codex Santa María Asunción, the so-called Uppsala Map, and the Maguey Plan (from Azcapotzalco). In alphabetic Nahuatl documentation, The Testaments of Culhuacan from the late sixteenth century have numerous references to chinampas as property that individuals bequeathed to their heirs in written wills. There are still remnants of the chinampa system in Xochimilco, the southern portion of greater Mexico City. Chinampas have been promoted as a model for modern sustainable agriculture, although some sources have disputed the applicability of this model. One anthropologist, for instance, reports that attempts by Mexico to develop chinampas among the Chontal Maya people in the 1970s failed until the technicians modified their goals in order to suit the Chontales' interests. Construction: According to Antonio Vera, through the UH Hilo website, within the framework of chinampas, there was two versions; inland and irrigated chinampas. Inland’s are created on banks, irrigated is built on water. Through steps, the structure of chinampas is to locate shallow land by the bank and surround said area with stakes of a common wetland tree [ahuejote]. The urbanization of Mexico lost this tradition and new challenges are created within the urbanization of Mexico. (https://hilo.hawaii.edu/nihopeku/2018/02/02/chinampa-an-ancient-agricultural-system/) Modern chinampas As of 1998, chinampas are still present in San Gregorio, a small town east of Xochimilco, in addition to San Luis, Tlahuac, and Mixquic. Although many of these gardens were constructed and thoroughly tended to from the Postclassic Period through the Spanish conquest, many of these plots of land still exist and are in active use. Many of these chinampas have been allowed by present-day farmers to become overgrown. Some choose to use canoes to farm, but many are becoming increasingly dependent on wheelbarrows and bicycles for transportation. Other fields, such as some located in San Gregorio and San Luis areas, have been deliberately filled up. As the canals dry up, several of the fields are naturally joined. Although not used for their original purpose, they are commonly used for cattle feed. Other fields, both dried and surrounded by canals, produce foods such as lettuce, cilantro, spinach, chard, squash, parsley, coriander, cauliflower, celery, mint, chives, rosemary, corn, and radishes. The young leaves of quelites and quintoniles, which are often mistaken for weeds, are grown and harvested as ingredients of sauces. Flowers also continue to be grown on these plots. Some chinampa fields are also used as tourist sites. Challenges Although many locals and farmers are happy to return to their agricultural roots, they are faced with several challenges. During the Spanish conquest, many lakes were drained, limiting their agricultural capacity, such as the lake at Xochimilco. In addition, in 1985 an earthquake struck, further damaging several canals. Other challenges include limited water supply, the use of pesticides, climate change, urban sprawl, and water pollution caused by untreated sewage and toxic waste.	Technology	Buildings and infrastructure	null
312293	https://en.wikipedia.org/wiki/Liouville%27s%20theorem%20%28complex%20analysis%29	Liouville's theorem (complex analysis)	In complex analysis, Liouville's theorem, named after Joseph Liouville (although the theorem was first proven by Cauchy in 1844), states that every bounded entire function must be constant. That is, every holomorphic function for which there exists a positive number such that for all is constant. Equivalently, non-constant holomorphic functions on have unbounded images. The theorem is considerably improved by Picard's little theorem, which says that every entire function whose image omits two or more complex numbers must be constant. Statement Liouville's theorem: Every holomorphic function for which there exists a positive number such that for all is constant. More succinctly, Liouville's theorem states that every bounded entire function must be constant. Proof This important theorem has several proofs. A standard analytical proof uses the fact that holomorphic functions are analytic. Another proof uses the mean value property of harmonic functions. The proof can be adapted to the case where the harmonic function is merely bounded above or below. See Harmonic function#Liouville's theorem. Corollaries Fundamental theorem of algebra There is a short proof of the fundamental theorem of algebra using Liouville's theorem. No entire function dominates another entire function A consequence of the theorem is that "genuinely different" entire functions cannot dominate each other, i.e. if and are entire, and everywhere, then for some complex number . Consider that for the theorem is trivial so we assume . Consider the function . It is enough to prove that can be extended to an entire function, in which case the result follows by Liouville's theorem. The holomorphy of is clear except at points in . But since is bounded and all the zeroes of are isolated, any singularities must be removable. Thus can be extended to an entire bounded function which by Liouville's theorem implies it is constant. If f is less than or equal to a scalar times its input, then it is linear Suppose that is entire and , for . We can apply Cauchy's integral formula; we have that where is the value of the remaining integral. This shows that is bounded and entire, so it must be constant, by Liouville's theorem. Integrating then shows that is affine and then, by referring back to the original inequality, we have that the constant term is zero. Non-constant elliptic functions cannot be defined on the complex plane The theorem can also be used to deduce that the domain of a non-constant elliptic function cannot be . Suppose it was. Then, if and are two periods of such that is not real, consider the parallelogram whose vertices are 0, , , and . Then the image of is equal to . Since is continuous and is compact, is also compact and, therefore, it is bounded. So, is constant. The fact that the domain of a non-constant elliptic function cannot be is what Liouville actually proved, in 1847, using the theory of elliptic functions. In fact, it was Cauchy who proved Liouville's theorem. Entire functions have dense images If is a non-constant entire function, then its image is dense in . This might seem to be a much stronger result than Liouville's theorem, but it is actually an easy corollary. If the image of is not dense, then there is a complex number and a real number such that the open disk centered at with radius has no element of the image of . Define Then is a bounded entire function, since for all , So, is constant, and therefore is constant. On compact Riemann surfaces Any holomorphic function on a compact Riemann surface is necessarily constant. Let be holomorphic on a compact Riemann surface . By compactness, there is a point where attains its maximum. Then we can find a chart from a neighborhood of to the unit disk such that is holomorphic on the unit disk and has a maximum at , so it is constant, by the maximum modulus principle. Remarks Let be the one-point compactification of the complex plane . In place of holomorphic functions defined on regions in , one can consider regions in . Viewed this way, the only possible singularity for entire functions, defined on , is the point . If an entire function is bounded in a neighborhood of , then is a removable singularity of , i.e. cannot blow up or behave erratically at . In light of the power series expansion, it is not surprising that Liouville's theorem holds. Similarly, if an entire function has a pole of order at —that is, it grows in magnitude comparably to in some neighborhood of —then is a polynomial. This extended version of Liouville's theorem can be more precisely stated: if for sufficiently large, then is a polynomial of degree at most . This can be proved as follows. Again take the Taylor series representation of , The argument used during the proof using Cauchy estimates shows that for all , So, if , then Therefore, . Liouville's theorem does not extend to the generalizations of complex numbers known as double numbers and dual numbers.	Mathematics	Complex analysis	null
312299	https://en.wikipedia.org/wiki/Avalanche%20photodiode	Avalanche photodiode	An avalanche photodiode (APD) is a highly sensitive type of photodiode, which in general are semiconductor diodes that convert light into electricity via interband excitation coupled with impact ionization. APDs use materials and a structure optimised for operating with high reverse bias, approaching the reverse breakdown voltage, such that charge carriers generated by the photoelectric effect are multiplied by an avalanche breakdown; thus they can be used to detect relatively small amounts of light. From a functional standpoint, they can be regarded as the semiconductor analog of photomultiplier tubes; unlike solar cells, they are not optimised for generating electricity from light but rather for detection of incoming photons. Typical applications for APDs are laser rangefinders, long-range fiber-optic telecommunication, positron emission tomography, and particle physics. History The avalanche photodiode was invented by Japanese engineer Jun-ichi Nishizawa in 1952. However, study of avalanche breakdown, micro-plasma defects in silicon and germanium and the investigation of optical detection using p-n junctions predate this patent. Principle of operation Photodiodes generally operate by impact ionization, whereby a photon provides the energy to separate charge carriers in the semiconductor material into a positive and negative pair, which can thus cause a charge flow through the diode. By applying a high reverse bias voltage, any photoelectric effect in the diode can be multiplied by the avalanche effect. Thus, the APD can be thought of as applying a high gain effect to the induced photocurrent. In general, the higher the reverse voltage, the higher the gain. A standard silicon APD typically can sustain 100–200 V of reverse bias before breakdown, leading to a gain factor of around 100. However, by employing alternative doping and bevelling (structural) techniques compared to traditional APDs, a it is possible to create designs where greater voltage can be applied (> 1500 V) before breakdown is reached, and hence a greater operating gain (> 1000) is achieved. Among the various expressions for the APD multiplication factor (M), an instructive expression is given by the formula where L is the space-charge boundary for electrons, and is the multiplication coefficient for electrons (and holes). This coefficient has a strong dependence on the applied electric field strength, temperature, and doping profile. Since APD gain varies strongly with the applied reverse bias and temperature, it is necessary to closely monitor the reverse voltage to keep a stable gain. Geiger mode counting If very high gain is needed (105 to 106), detectors related to APDs called SPADs (single-photon avalanche diodes) can be used and operated with a reverse voltage above a typical APD's breakdown voltage. In this case, the photodetector needs to have its signal current limited and quickly diminished. Active and passive current-quenching techniques have been used for this purpose. SPADs that operate in this high-gain regime are sometimes referred to being in Geiger mode. This mode is particularly useful for single-photon detection, provided that the dark count event rate and afterpulsing probability are sufficiently low. Materials In principle, any semiconductor material can be used as a multiplication region: Silicon will detect in the visible and near infrared, with low multiplication noise (excess noise). Germanium (Ge) will detect infrared out to a wavelength of 1.7 μm, but has high multiplication noise. InGaAs will detect out to longer than 1.6 μm and has less multiplication noise than Ge. It is normally used as the absorption region of a heterostructure diode, most typically involving InP as a substrate and as a multiplication layer. This material system is compatible with an absorption window of roughly 0.9–1.7 μm. InGaAs exhibits a high absorption coefficient at the wavelengths appropriate to high-speed telecommunications using optical fibers, so only a few micrometres of InGaAs are required for nearly 100% light absorption. The excess noise factor is low enough to permit a gain-bandwidth product in excess of 100 GHz for a simple InP/InGaAs system, and up to 400 GHz for InGaAs on silicon. Therefore, high-speed operation is possible: commercial devices are available to speeds of at least 10 Gbit/s. Gallium-nitride–based diodes have been used for operation with ultraviolet light. HgCdTe-based diodes operate in the infrared, typically at wavelengths up to about 14 μm, but require cooling to reduce dark currents. Very low excess noise can be achieved in this material system. Structure APDs are often not constructed as simple p-n junctions but have more complex designs such as p+-i-p-n+. Performance limits APD applicability and usefulness depends on many parameters. Two of the larger factors are: quantum efficiency, which indicates how well incident optical photons are absorbed and then used to generate primary charge carriers; and total leakage current, which is the sum of the dark current, photocurrent and noise. Electronic dark-noise components are series and parallel noise. Series noise, which is the effect of shot noise, is basically proportional to the APD capacitance, while the parallel noise is associated with the fluctuations of the APD bulk and surface dark currents. Gain noise, excess noise factor Another noise source is the excess noise factor, ENF. It is a multiplicative correction applied to the noise that describes the increase in the statistical noise, specifically Poisson noise, due to the multiplication process. The ENF is defined for any device, such as photomultiplier tubes, silicon solid-state photomultipliers, and APDs, that multiplies a signal, and is sometimes referred to as "gain noise". At a gain M, it is denoted by ENF(M) and can often be expressed as where is the ratio of the hole impact ionization rate to that of electrons. For an electron multiplication device it is given by the hole impact ionization rate divided by the electron impact ionization rate. It is desirable to have a large asymmetry between these rates to minimize ENF(M), since ENF(M) is one of the main factors that limit, among other things, the best possible energy resolution obtainable. Conversion noise, Fano factor The noise term for an APD may also contain a Fano factor, which is a multiplicative correction applied to the Poisson noise associated with the conversion of the energy deposited by a charged particle to the electron-hole pairs, which is the signal before multiplication. The correction factor describes the decrease in the noise, relative to Poisson statistics, due to the uniformity of conversion process and the absence of, or weak coupling to, bath states in the conversion process. In other words, an "ideal" semiconductor would convert the energy of the charged particle into an exact and reproducible number of electron hole pairs to conserve energy; in reality, however, the energy deposited by the charged particle is divided into the generation of electron hole pairs, the generation of sound, the generation of heat, and the generation of damage or displacement. The existence of these other channels introduces a stochastic process, where the amount of energy deposited into any single process varies from event to event, even if the amount of energy deposited is the same. Further influences The underlying physics associated with the excess noise factor (gain noise) and the Fano factor (conversion noise) is very different. However, the application of these factors as multiplicative corrections to the expected Poisson noise is similar. In addition to excess noise, there are limits to device performance associated with the capacitance, transit times and avalanche multiplication time. The capacitance increases with increasing device area and decreasing thickness. The transit times (both electrons and holes) increase with increasing thickness, implying a tradeoff between capacitance and transit time for performance. The avalanche multiplication time times the gain is given to first order by the gain-bandwidth product, which is a function of the device structure and most especially .	Technology	Components	null
312335	https://en.wikipedia.org/wiki/Tarantula%20Nebula	Tarantula Nebula	The Tarantula Nebula (also known as 30 Doradus) is a large H II region in the Large Magellanic Cloud (LMC), forming its south-east corner (from Earth's perspective). Discovery The Tarantula Nebula was observed by Nicolas-Louis de Lacaille during an expedition to the Cape of Good Hope between 1751 and 1753. He cataloged it as the second of the "Nebulae of the First Class", "Nebulosities not accompanied by any star visible in the telescope of two feet". It was described as a diffuse nebula 20' across. Johann Bode included the Tarantula in his 1801 Uranographia star atlas and listed it in the accompanying Allgemeine Beschreibung und Nachweisung der Gestirne catalog as number 30 in the constellation "Xiphias or Dorado". Instead of being given a stellar magnitude, it was noted to be nebulous. The name Tarantula Nebula arose in the mid-20th century from its appearance in deep photographic exposures. 30 Doradus has often been treated as the designation of a star, or of the central star cluster NGC 2070, but is now generally treated as referring to the whole nebula area of the Tarantula Nebula. Properties The Tarantula Nebula has an apparent magnitude of 8. Considering its distance of about 49 kpc (160,000 light-years), this is an extremely luminous non-stellar object. Its luminosity is so great that if it were as close to Earth as the Orion Nebula, the Tarantula Nebula would cast visible shadows. In fact, it is the most active starburst region known in the Local Group of galaxies. It is also one of the largest H II regions in the Local Group with an estimated diameter around 200 to 570 pc (650 to 1860 light years), and also because of its very large size, it is sometimes described as the largest. However, other H II regions such as NGC 604, which is in the Triangulum Galaxy, could be larger. The nebula resides on the leading edge of the LMC where ram pressure stripping, and the compression of the interstellar medium likely resulting from this, is at a maximum. NGC 2070 30 Doradus has at its centre the star cluster NGC 2070 which includes the compact concentration of stars known as R136 that produces most of the energy that makes the nebula visible. The estimated mass of the cluster is 450,000 solar masses, suggesting it will likely become a globular cluster in the future. In addition to NGC 2070, the Tarantula Nebula contains several other star clusters including the much older Hodge 301. The most massive stars of Hodge 301 have already exploded in supernovae. Supernova 1987A The closest supernova observed since the invention of the telescope, Supernova 1987A, occurred in the outskirts of the Tarantula Nebula. There is a prominent supernova remnant enclosing the open cluster NGC 2060. Still, the remnants of many other supernovae are difficult to detect in the complex nebulosity. Black hole VFTS 243 An x-ray quiet black hole was discovered in the Tarantula Nebula, the first outside of the Milky Way Galaxy that does not radiate strongly. The black hole has a mass of at least 9 solar masses and is in a circular orbit with its 25 solar mass blue giant companion VFTS 243.	Physical sciences	Notable nebulae	Astronomy
312461	https://en.wikipedia.org/wiki/Unit%20of%20account	Unit of account	In economics, unit of account is one of the functions of money. A unit of account is a standard numerical monetary unit of measurement of the market value of goods, services, and other transactions. Also known as a "measure" or "standard" of relative worth and deferred payment, a unit of account is a necessary prerequisite for the formulation of commercial agreements that involve debt. Money acts as a standard measure and a common denomination of trade. It is thus a basis for quoting and bargaining of prices. It is necessary for developing efficient accounting systems. Economics Unit of account in economics allows a somewhat meaningful interpretation of prices, costs, and profits, so that an entity can monitor its own performance. It allows shareholders to make sense of its past performance and have an idea of its future profitability. The use of money, as a relatively stable unit of measure, can tend to drive market economies toward efficiency. Historically, prices were often given in a dominant currency used as a unit of account, but transactions actually settled by using a variety of coins that were available, and often goods, all converted into their value in the unit of account. Many international transactions continue to be settled in this way, using a national value (most often expressed in the US dollar or euro) but with the actual settlement in something else. In historical cost accounting, currencies are assumed to be perfectly stable in real value during non-hyperinflationary conditions under in terms of which the stable measuring unit assumption is applied. The Daily Consumer Price Index (Daily CPI) – or a monetized daily indexed unit of account – can be used to index monetary values on a daily basis when it is required to maintain the purchasing power or real value of monetary values constant during inflation and deflation. Problems Money is rarely perfectly stable in real value which is the fundamental problem with traditional historical cost accounting which is based on the stable measuring unit assumption. The unit of account in economics suffers from the pitfall of not being stable in real value over time because money is generally not perfectly stable in real value during inflation and deflation. Inflation destroys the assumption that the real value of the unit of account is stable which is the basis of classic accountancy. In such circumstances, historical values registered in accountancy books become heterogeneous amounts measured in different units. The use of such data under traditional accounting methods without previous correction can lead to confusing — (or even meaningless) — results. History Historic examples of units of measure include the livre tournois, used in France from 1302 to 1794 whether or not livre coins were minted. In the 14th century Naples used the grossi gigliati, and Bohemia used the Prague groschen. (2021) At any one time there might be two or three units of account in one region based on the local base, silver and sometimes gold coins, and each often expressed in L.S.D units in ratio 240:12:1. The Florentine gold florin, the French franc and the electoral rheingulden all became pounds (240 denari) of account. Units of account would often survive over 100 years despite the original coins changing composition and availability (e.g. the Castilian maravedi). In 1921, Henry Ford proposed the use of energy as the basis for currency instead of Gold Standard. Thomas Edison similarly put forward commodities as a basis. At the onset of the Great Depression, John P. Norton restated the "Electric Dollars" standard alongside gold. A modern unit of account is the European Currency Unit, used in the European Union from 1979 to 1998; its replacement in 1999, the Euro, was also just a unit of account until the introduction of notes and coins in 2002. Unit of account is the main way of calculating a carrier or ship owner's liability in relation to carriage of goods contracts in which the Hague-Visby Rules apply. In economics, a standard unit of account is used for statistical purposes to describe economic activity. Indexes such as GDP and the CPI are so broad in their scope that compiling them would be impossible without a standard unit of account. After being compiled, these figures are often used to guide governmental policy; especially monetary and fiscal policy. In calculating the opportunity cost of a policy, a standard unit of account allows for the creation of a composite good. A composite good is a theoretical abstraction that represents an aggregation of all other opportunities that are not realized by the first good. It allows an economic decision's benefits to be weighed against the costs of all other possible goods in that society, without having to refer to any directly. Often, this is most easily accomplished with money. Finance The use of a unit of account in financial accounting, according to the American business model, allows investors to invest capital into those companies that provide the highest rate of return. The use of a unit of account in managerial accounting enables firms to choose between activities that yield the highest profit. Accounting The unit of account in financial accounting refers to the words used to describe the specific assets and liabilities that are reported in financial statements rather than the units used to measure them. That is, unit of account refers to the object of recognition or display whereas unit of measure refers to the tool for measuring it. Unit of measure and unit of account are sometimes treated as synonyms in financial accounting and economics. Unit of measure in financial accounting refers to the monetary unit to be used; that is, whether it should be nominal units of money as opposed to units that are adjusted for changes in purchasing power over time.	Physical sciences	Measurement systems	Basics and measurement
312678	https://en.wikipedia.org/wiki/Footpath	Footpath	A footpath (also pedestrian way, walking trail, nature trail) is a type of thoroughfare that is intended for use only by pedestrians and not other forms of traffic such as motorized vehicles, bicycles and horses. They can be found in a wide variety of places, from the centre of cities, to farmland, to mountain ridges. Urban footpaths are usually paved, may have steps, and can be called alleys, lanes, steps, etc. National parks, nature preserves, conservation areas and other protected wilderness areas may have footpaths (trails) that are restricted to pedestrians. The term footpath can also describe a pavement/sidewalk in some English-speaking countries (such as Australia, New Zealand, and Ireland). A footpath can also take the form of a footbridge, linking two places across a river. Origins and history Public footpaths are rights of way originally created by people walking across the land to work, market, the next village, church, and school. This includes mass paths and corpse roads. Some footpaths were also created by those undertaking a pilgrimage. Examples of the latter are the Pilgrim's Way in England and Pilgrim's Route (St. Olav's Way or the Old Kings' Road) in Norway. Some landowners allow access over their land without dedicating a right of way. These permissive paths are often indistinguishable from normal paths, but they are usually subject to restrictions. Such paths are often closed at least once a year, so that a permanent right of way cannot be established in law. A mass path is a pedestrian track or road connecting destinations frequently used by rural communities, most usually the destination of Sunday Mass. They were most common during the centuries that preceded motorised transportation in Western Europe, and in particular the British Isles and the Netherlands (where such a path is called "kerkenpad" (lit. Church path). Mass paths typically included stretches crossing fields of neighboring farmers and were likely to contain stiles, when crossing fences or other boundaries, or plank footbridges to cross ditches. Some mass paths are still used today in the Republic of Ireland, but are usually subject to Ireland's complicated rights of way law. Corpse roads provided a practical means for transporting corpses, often from remote communities, to cemeteries that had burial rights, such as parish churches and chapels of ease. In Great Britain, such routes can also be known by a number of other names: bier road, burial road, coffin road, coffin line, lyke or lych way, funeral road, procession way, corpse way, etc. Nowadays footpaths are mainly used for recreation and have been frequently linked together, along with bridle paths and newly created footpaths, to create long-distance trails. Also, organizations have been formed in various countries to protect the right to use public footpaths, including the Ramblers Association and the Open Spaces Society in England. Footpaths are now also found in botanic gardens, arboretums, regional parks, conservation areas, wildlife gardens, and open-air museums. There are also educational trails, themed walks, sculpture trails and historic interpretive trails. Rights of way In England and Wales, public footpaths are rights of way on which pedestrians have a legally protected right to travel. Other public rights of way in England and Wales, such as bridleways, byways, towpaths, and green lanes are also used by pedestrians. In Scotland there is no legal distinction between a footpath and a bridleway and it is generally accepted that cyclists and horse riders may follow any right of way with a suitable surface. The law is different in both Northern Ireland and the Republic of Ireland and there are far fewer rights of way in Ireland as a whole (see Keep Ireland Open). Definitive path maps Footpaths and other rights of way in England and Wales are shown on definitive maps. A definitive map is a record of public rights of way in England and Wales. In law it is the definitive record of where a right of way is located. The highway authority (normally the county council, or unitary authority in areas with a one-tier system) has a statutory duty to maintain a definitive map, though in national parks the national park authority usually maintains the map. The Inner London boroughs are exempt from the statutory duty though they have the powers to maintain a map: currently none does so. Currently, the number of footpaths in the UK totals 427,301 (around 81% of all rights of way) with a net combined route length of 105,125 miles. In Scotland different legislation applies and there is no legally recognised record of rights of way. However, there is a National Catalogue of Rights of Way (CROW), compiled by the Scottish Rights of Way and Access Society (Scotways), in partnership with Scottish Natural Heritage, and the help of local authorities. Open Spaces Society The Open Spaces Society is a charitable British organisation that works to protect public rights of way and open spaces in the United Kingdom, such as common land and village greens. It is Britain's oldest national conservation body. The society was founded as the Commons Preservation Society and merged with the National Footpaths Society in 1899, and adopted their present name. Much of the Open Spaces Society's work is concerned with the preservation and creation of public paths. Before the introduction of definitive maps of public paths in the early 1950s, the public did not know where paths were, and the Open Spaces Society helped the successful campaign for paths to be shown on Ordnance Survey maps. It advises the Department for Environment, Food and Rural Affairs and National Assembly for Wales on applications for works on common land. Local authorities are legally required to consult the society whenever there is a proposal to alter the route of a public right of way. The Ramblers are another British organisation concerned with the protection of footpaths. Urban footpaths There are a variety of footpaths in urban settings, including paths along streams and rivers, through parks and across commons. Another type is the alley, normally providing access to the rear of properties or connecting built-up roads not easily reached by vehicles. Towpaths are another kind of urban footpath, but they are often shared with cyclists. A typical footpath in a park is found along the seawall in Stanley Park, Vancouver, British Columbia, Canada. This is a segregated path, with one lane for skaters and cyclists and the other for pedestrians. In the US and Canada, where urban sprawl has begun to strike even the most rural communities, developers and local leaders are currently striving to make their communities more conducive to non-motorized transportation through the use of less traditional paths. The Robert Wood Johnson Foundation has established the Active Living by Design program to improve the livability of communities in part through developing trails, The Upper Valley Trails Alliance has done similar work on traditional trails, while the Somerville Community Path and related paths, are examples of urban initiatives. In St. John's, Newfoundland, Canada The Grand Concourse, is an integrated walkway system that has over of footpaths which link every major park, river, pond, and green space in six municipalities. In London, England, there are several long-distance walking routes which combine footpaths and roads to link green spaces. These include the Capital Ring, London Outer Orbital Path and the Jubilee Walkway, the use of which have been endorsed by Transport for London. Alley and steps An alley is a narrow, usually paved, pedestrian path, often between the walls of buildings in towns and cities. This type is usually short and straight, and on steep ground can consist partially or entirely of steps. In older cities and towns in Europe, alleys are often what is left of a medieval street network, or a right of way or ancient footpath. Similar paths also exist in some older North American towns and cities. In some older urban development in North America lanes at the rear of houses, to allow for deliveries and garbage collection, are called alleys. Alleys may be paved, or unpaved, and a blind alley is a cul-de-sac. Some alleys are roofed because they are within buildings, such as the traboules of Lyon, or when they are a pedestrian passage through railway embankments in Britain. The latter follow the line of rights-of way that existed before the railway was built. Because of topography, steps (stairs) are the predominant form of alley in hilly cities and towns. This includes Pittsburgh (see Steps of Pittsburgh), Cincinnati (see Steps of Cincinnati), Portland, Oregon, Seattle, and San Francisco in the United States, as well as Hong Kong, and Rome. Long-distance paths Footpaths (and other rights of way) have been combined, and new paths created, so as to produce long-distance walking routes in a number of countries. These can be rural in nature, such as the Essex Way, in southern England, which crosses farmland, or urban as with various routes in London, England, or along a coastline like the South West Coast Path in the West of England, or in the high mountains, like the Pacific Crest Trail in the US, which reaches at Forester Pass in the Sierra Nevada. Maintenance Many footpaths require some maintenance. Most rural paths have an earth or grass surface with stiles, and or gates, including kissing gates. A few will have stepping stones, fords, or bridges. Urban footpaths may be constructed of masonry, brick, concrete, asphalt, cut stone or wood boardwalk. Crushed rock, decomposed granite, fine wood chips are also used. The construction materials can vary over the length of the footpath and may start with a well constructed hard surface in an urban area, and end with an inexpensive soft or loose surface in the countryside. Stairs or steps are sometimes found in urban alleys, or cliff paths to beaches. Issues The main issues in urban areas include maintenance, litter, crime, and lighting after dark. In the countryside there are issues relating to conflicts between walkers and livestock, and these occasionally result in people being injured or even killed. Dogs often contribute to such conflicts – see in England and Wales The Dogs (Protection of Livestock) Act 1953. Also footpaths in remote locations can be difficult to maintain and a route along a country path can be impeded by ploughing, crops, overgrown vegetation, illegal barriers (including barbed wire), damaged stiles, etc. Confrontation with landowners in the UK There have been numerous problems over the years in England and Wales with landowners. One notable example was with the millionaire property tycoon Nicholas Van Hoogstraten who had a long-standing dislike of and dispute with ramblers, describing them as "scum of the earth". In 1999 Hoogstraten erected a large fence across a footpath on his country estate in East Sussex. Local ramblers staged a protest against the erection of the fence outside the boundary of Van Hoogstraten's estate. On 10 February 2003 and after a 13-year battle and numerous legal proceedings, the path was finally re-opened. Isle of Man Another conflict involved Jeremy Clarkson, a TV presenter and Top Gear host who lives on the Isle of Man. He became frustrated at the lack of privacy at his home when ramblers deviated from a pathway to take photographs of his dwelling. Clarkson's property bordered a small 250-metre strip of land that had no definitive status as a public right of way but was used by walkers regardless. Clarkson aimed to close access to this small strip of his land, thereby forcing ramblers to take a small diversion to stick to the official public right of way and therefore protecting his claimed right to privacy on his own property. In May 2010 the former transport minister, Hon. David Anderson MHK, accepted the conclusions of a public inquiry that all except five of the paths claimed at the inquiry as public rights of way have been dedicated as public rights of way and should be added to the definitive map.	Technology	Road infrastructure	null
312833	https://en.wikipedia.org/wiki/Dots%20per%20inch	Dots per inch	Dots per inch (DPI, or dpi) is a measure of spatial printing, video or image scanner dot density, in particular the number of individual dots that can be placed in a line within the span of . Similarly, dots per centimetre (d/cm or dpcm) refers to the number of individual dots that can be placed within a line of . DPI measurement in printing DPI is used to describe the resolution number of dots per inch in a digital print and the printing resolution of a hard copy print dot gain, which is the increase in the size of the halftone dots during printing. This is caused by the spreading of ink on the surface of the media. Up to a point, printers with higher DPI produce clearer and more detailed output. A printer does not necessarily have a single DPI measurement; it is dependent on print mode, which is usually influenced by driver settings. The range of DPI supported by a printer is most dependent on the print head technology it uses. A dot matrix printer, for example, applies ink via tiny rods striking an ink ribbon, and has a relatively low resolution, typically in the range of . An inkjet printer sprays ink through tiny nozzles, and is typically capable of 300–720 DPI. A laser printer applies toner through a controlled electrostatic charge, and may be in the range of 600 to 2,400 DPI. The DPI measurement of a printer often needs to be considerably higher than the pixels per inch (PPI) measurement of a video display in order to produce similar-quality output. This is due to the limited range of colours for each dot typically available on a printer. At each dot position, the simplest type of color printer can either print no dot, or print a dot consisting of a fixed volume of ink in each of four color channels (typically CMYK with cyan, magenta, yellow and black ink) or 24 = 16 colours on laser, wax and most inkjet printers, of which only 14 or 15 (or as few as 8 or 9) may be actually discernible depending on the strength of the black component, the strategy used for overlaying and combining it with the other colours, and whether it is in "color" mode. Higher-end inkjet printers can offer 5, 6 or 7 ink colours giving 32, 64 or 128 possible tones per dot location (and again, it can be that not all combinations will produce a unique result). Contrast this to a standard sRGB monitor where each pixel produces 256 intensities of light in each of three channels (RGB). While some color printers can produce variable drop volumes at each dot position, and may use additional ink-color channels, the number of colours is still typically less than on a monitor. Most printers must therefore produce additional colours through a halftone or dithering process, and rely on their base resolution being high enough to "fool" the human observer's eye into perceiving a patch of a single smooth colour. The exception to this rule is dye-sublimation printers, which can apply a much more variable amount of dye—close to or exceeding the number of the 256 levels per channel available on a typical monitor—to each "pixel" on the page without dithering, but with other limitations: lower spatial resolution (typically 200 to 300 dpi), which can make text and lines look somewhat rough lower output speed (a single page requiring three or four complete passes, one for each dye colour, each of which may take more than fifteen seconds—generally quicker, however, than most inkjet printers' "photo" modes) a wasteful (and, for confidential documents, insecure) dye-film roll cartridge system occasional color registration errors (mainly along the long axis of the page), which necessitate recalibrating the printer to account for slippage and drift in the paper feed system. These disadvantages mean that, despite their marked superiority in producing good photographic and non-linear diagrammatic output, dye-sublimation printers remain niche products, and thus other devices using higher resolution, lower color depth, and dither patterns remain the norm. This dithered printing process could require a region of four to six dots (measured across each side) to accurately reproduce the color in a single pixel. An image that is 100 pixels wide may need to be 400 to 600 dots in width in the printed output; if a 100 × 100-pixel image is to be printed in a one-inch square, the printer must be capable of 400 to 600 dots per inch to reproduce the image. As such, 600 dpi (sometimes 720) is now the typical output resolution of entry-level laser printers and some utility inkjet printers, with 1,200–1,440 and 2,400–2,880 being common "high" resolutions. This contrasts with the 300–360 (or 240) dpi of early models, and the approximate 200 dpi of dot-matrix printers and fax machines, which gave faxed and computer-printed documents—especially those that made heavy use of graphics or coloured block text—a characteristic "digitized" appearance, because of their coarse, obvious dither patterns, inaccurate colours, loss of clarity in photographs, and jagged ("aliased") edges on some text and line art. DPI or PPI in digital image files In printing, DPI (dots per inch) refers to the output resolution of a printer or imagesetter, and PPI (pixels per inch) refers to the input resolution of a photograph or image. DPI refers to the physical dot density of an image when it is reproduced as a real physical entity, for example printed onto paper. A digitally stored image has no inherent physical dimensions, measured in inches or centimetres. Some digital file formats record a DPI value, or more commonly a PPI (pixels per inch) value, which is to be used when printing the image. This number lets the printer or software know the intended size of the image, or in the case of scanned images, the size of the original scanned object. For example, a bitmap image may measure 1,000 × 1,000 pixels, a resolution of 1 megapixel. If it is labelled as 250 PPI, that is an instruction to the printer to print it at a size of 4 × 4 inches. Changing the PPI to 100 in an image editing program would tell the printer to print it at a size of 10 × 10 inches. However, changing the PPI value would not change the size of the image in pixels which would still be 1,000 × 1,000. An image may also be resampled to change the number of pixels and therefore the size or resolution of the image, but this is quite different from simply setting a new PPI for the file. For vector images, since the file is resolution independent, there is no need to resample the image before resizing it as it prints equally well at all sizes. However, there is still a target printing size. Some image formats, such as Photoshop format, can contain both bitmap and vector data in the same file. Adjusting the PPI in a Photoshop file will change the intended printing size of the bitmap portion of the data and also change the intended printing size of the vector data to match. This way the vector and bitmap data maintain a consistent size relationship when the target printing size is changed. Text stored as outline fonts in bitmap image formats is handled in the same way. Other formats, such as PDF, are primarily vector formats that can contain images, potentially at a mixture of resolutions. In these formats the target PPI of the bitmaps is adjusted to match when the target print size of the file is changed. This is the converse of how it works in a primarily bitmap format like Photoshop, but has exactly the same result of maintaining the relationship between the vector and bitmap portions of the data. Computer monitor DPI standards Since the 1980s, Macs have set the default display "DPI" to 72 PPI, while the Microsoft Windows operating system has used a default of 96 PPI. These default specifications arose out of the problems rendering standard fonts in the early display systems of the 1980s, including the IBM-based CGA, EGA, VGA and 8514 displays as well as the Macintosh displays featured in the 128K computer and its successors. The choice of 72 PPI by Macintosh for their displays arose from existing convention: the official 72 points per inch mirrored the 72 pixels per inch that appeared on their display screens. (Points are a physical unit of measure in typography, dating from the days of printing presses, where 1 point by the modern definition is of the international inch (25.4 mm), which therefore makes 1 point approximately 0.0139 in or 352.8 μm). Thus, the 72 pixels per inch seen on the display had exactly the same physical dimensions as the 72 points per inch later seen on a printout, with 1 pt in printed text equal to 1 px on the display screen. As it is, the Macintosh 128K featured a screen measuring 512 pixels in width by 342 pixels in height, and this corresponded to the width of standard office paper (512 px ÷ 72 px/in ≈ 7.1 in, with a 0.7 in margin down each side when assuming in × 11 in North American paper size; in the rest of the world, it is 210 mm × 297 mm – called A4. B5 is 176 mm × 250 mm). A consequence of Apple's decision was that the widely used 10-point fonts from the typewriter era had to be allotted 10 display pixels in em height, and 5 display pixels in x-height. This is technically described as 10 pixels per em (PPEm). This made 10-point fonts be rendered crudely and made them difficult to read on the display screen, particularly the lowercase characters. Furthermore, there was the consideration that computer screens are typically viewed (at a desk) at a distance 30% greater than printed materials, causing a mismatch between the perceived sizes seen on the computer screen and those on the printouts. Microsoft tried to solve both problems with a hack that has had long-term consequences for the understanding of what DPI and PPI mean. Microsoft began writing its software to treat the screen as though it provided a PPI characteristic that is of what the screen actually displayed. Because most screens at the time provided around 72 PPI, Microsoft essentially wrote its software to assume that every screen provides 96 PPI (because 72 × = 96). The short-term gain of this trickery was twofold: It would seem to the software that one-third more pixels were available for rendering an image, thereby allowing for bitmap fonts to be created with greater detail. On every screen that actually provided 72 PPI, each graphical element (such as a character of text) would be rendered at a size one third larger than it "should" be, thereby allowing a person to sit a comfortable distance from the screen. However, larger graphical elements meant less screen space was available for programs to draw. Indeed, the default 720-pixel wide mode of a Hercules mono graphics adaptor (the one-time gold standard for high resolution PC graphics) – or a "tweaked" VGA adaptor – provided an apparent -inch page width at this resolution. However, the more common and colour-capable display adaptors of the time all provided a 640-pixel wide image in their high resolution modes, enough for a bare inches at 100% zoom, with barely any greater visible page height – a maximum of 5 inches, versus . Consequently, the default margins in Microsoft Word were set, and still remain at 1 full inch on all sides of the page, keeping the "text width" for standard size printer paper within visible limits; despite most computer monitors now being both larger and finer-pitched, and printer paper transports having become more sophisticated, the Mac-standard half-inch borders remain listed in Word 2010's page layout presets as the "narrow" option (versus the 1-inch default). Without using supplemental, software-provided zoom levels, the 1:1 relationship between display and print size was (deliberately) lost; the availability of different-sized, user-adjustable monitors and display adaptors with varying output resolutions exacerbated this, as it was not possible to rely on a properly-adjusted "standard" monitor and adaptor having a known PPI. For example, a 12-inch Hercules monitor and adaptor with a thick bezel and a little underscan may offer 90 "physical" PPI, with the displayed image appearing nearly identical to hardcopy (assuming the H-scan density was properly adjusted to give square pixels) but a thin-bezel 14-inch VGA monitor adjusted to give a borderless display may be closer to 60, with the same bitmap image thus appearing 50% larger; yet, someone with an 8514 ("XGA") adaptor and the same monitor could achieve 100 DPI using its 1024-pixel wide mode and adjusting the image to be underscanned. A user who wanted to directly compare on-screen elements against those on an existing printed page by holding it up against the monitor would therefore first need to determine the correct zoom level to use, largely by trial and error, and often not be able to obtain an exact match in programs that only allowed integer per cent settings, or even fixed pre-programmed zoom levels. For the examples above, they may need to use respectively 94% (precisely, 93.75) – or , 63% (62.5) – or ; and 104% (104.167) – or , with the more commonly accessible 110% actually being a less precise match. Thus, for example, a 10-point font on a Macintosh (at 72 PPI) was represented with 10 pixels (i.e., 10 PPEm), whereas a 10-point font on a Windows platform (at 96 PPI) at the same zoom level is represented with 13 pixels (i.e., Microsoft rounded to 13 pixels, or 13 PPEm) – and, on a typical consumer grade monitor, would have physically appeared around to inch high instead of . Likewise, a 12-point font was represented with 12 pixels on a Macintosh, and 16 pixels (or a physical display height of maybe inch) on a Windows platform at the same zoom, and so on. The negative consequence of this standard is that with 96 PPI displays, there is no longer a one-to-one relationship between the font size in pixels and the printout size in points. This difference is accentuated on more recent displays that feature higher pixel densities. This has been less of a problem with the advent of vector graphics and fonts being used in place of bitmap graphics and fonts. Moreover, many Windows software programs have been written since the 1980s which assume that the screen provides 96 PPI. Accordingly, these programs do not display properly at common alternative resolutions such as 72 PPI or 120 PPI. The solution has been to introduce two concepts: logical PPI: The PPI that software claims a screen provides. This can be thought of as the PPI provided by a virtual screen created by the operating system. physical PPI: The PPI that a physical screen actually provides. Software programs render images to the virtual screen and then the operating system renders the virtual screen onto the physical screen. With a logical PPI of 96 PPI, older programs can still run properly regardless of the actual physical PPI of the display screen, although they may exhibit some visual distortion thanks to the effective 133.3% pixel zoom level (requiring either that every third pixel be doubled in width/height, or heavy-handed smoothing be employed). How Microsoft Windows handles DPI scaling Displays with high pixel densities were not common up to the Windows XP era. High DPI displays became mainstream around the time Windows 8 was released. Display scaling by entering a custom DPI irrespective of the display resolution has been a feature of Microsoft Windows since Windows 95. Windows XP introduced the GDI+ library which allows resolution-independent text scaling. In Microsoft Windows, the DPI higher than 96 DPI is called High DPI. Windows Vista introduced support for programs to declare themselves to the OS that they are high-DPI aware via a manifest file or using an API. For programs that do not declare themselves as DPI-aware, Windows Vista supports a compatibility feature called DPI virtualization so system metrics and UI elements are presented to applications as if they are running at 96 DPI and the Desktop Window Manager then scales the resulting application window to match the DPI setting. Windows Vista retains the Windows XP style scaling option which when enabled turns off DPI virtualization for all applications globally. DPI virtualization is a compatibility option as application developers are all expected to update their apps to support high DPI without relying on DPI virtualization. Windows Vista also introduces Windows Presentation Foundation. WPF .NET applications are vector-based, not pixel-based and are designed to be resolution-independent. Developers using the old GDI API and Windows Forms on .NET Framework runtime need to update their apps to be DPI aware and flag their applications as DPI-aware. Windows 7 adds the ability to change the DPI by doing only a log off, not a full reboot and makes it a per-user setting. Additionally, Windows 7 reads the pixel density related information from the EDID and automatically sets the system DPI value to match the monitor's physical pixel density, unless the effective resolution is less than 1024 × 768. Also, Windows 7 adds DirectWrite that optimised for monitors that larger than 1080p. In Windows 8, only the DPI scaling percentage is shown in the DPI changing dialog and the display of the raw DPI value has been removed. In Windows 8.1, the global setting to disable DPI virtualization (only use XP-style scaling) is removed and a per-app setting added for the user to disable DPI virtualization from the Compatibility tab. When the DPI scaling setting is set to be higher than 120 PPI (125%), DPI virtualization is enabled for all applications unless the application opts out of it by specifying a DPI aware flag (manifest) as "true" inside the EXE. Windows 8.1 retains a per-application option to disable DPI virtualization of an app. Windows 8.1 also adds the ability for different displays to use independent DPI scaling factors, although it calculates this automatically for each display and turns on DPI virtualization for all monitors at any scaling level. Windows 10 adds manual control over DPI scaling for individual monitors. Proposed metrication There are some ongoing efforts to abandon the DPI Image resolution unit in favour of a metric unit, giving the inter-dot spacing in dots per centimetre (px/cm or dpcm), as used in CSS3 media queries or micrometres (μm) between dots. A resolution of 72 DPI, for example, equals a resolution of about 28 dpcm or an inter-dot spacing of about 353 μm.	Physical sciences	Visual	Basics and measurement
312881	https://en.wikipedia.org/wiki/Action%20%28physics%29	Action (physics)	In physics, action is a scalar quantity that describes how the balance of kinetic versus potential energy of a physical system changes with trajectory. Action is significant because it is an input to the principle of stationary action, an approach to classical mechanics that is simpler for multiple objects. Action and the variational principle are used in Feynman's formulation of quantum mechanics and in general relativity. For systems with small values of action similar to the Planck constant, quantum effects are significant. In the simple case of a single particle moving with a constant velocity (thereby undergoing uniform linear motion), the action is the momentum of the particle times the distance it moves, added up along its path; equivalently, action is the difference between the particle's kinetic energy and its potential energy, times the duration for which it has that amount of energy. More formally, action is a mathematical functional which takes the trajectory (also called path or history) of the system as its argument and has a real number as its result. Generally, the action takes different values for different paths. Action has dimensions of energy × time or momentum × length, and its SI unit is joule-second (like the Planck constant h). Introduction Introductory physics often begins with Newton's laws of motion, relating force and motion; action is part of a completely equivalent alternative approach with practical and educational advantages. However, the concept took many decades to supplant Newtonian approaches and remains a challenge to introduce to students. Simple example For a trajectory of a ball moving in the air on Earth the action is defined between two points in time, and as the kinetic energy (KE) minus the potential energy (PE), integrated over time. The action balances kinetic against potential energy. The kinetic energy of a ball of mass is where is the velocity of the ball; the potential energy is where is the gravitational constant. Then the action between and is The action value depends upon the trajectory taken by the ball through and . This makes the action an input to the powerful stationary-action principle for classical and for quantum mechanics. Newton's equations of motion for the ball can be derived from the action using the stationary-action principle, but the advantages of action-based mechanics only begin to appear in cases where Newton's laws are difficult to apply. Replace the ball with an electron: classical mechanics fails but stationary action continues to work. The energy difference in the simple action definition, kinetic minus potential energy, is generalized and called the Lagrangian for more complex cases. Planck's quantum of action The Planck constant, written as or when including a factor of , is called the quantum of action. Like action, this constant has unit of energy times time. It figures in all significant quantum equations, like the uncertainty principle and the de Broglie wavelength. Whenever the value of the action approaches the Planck constant, quantum effects are significant. History Pierre Louis Maupertuis and Leonhard Euler working in the 1740s developed early versions of the action principle. Joseph Louis Lagrange clarified the mathematics when he invented the calculus of variations. William Rowan Hamilton made the next big breakthrough, formulating Hamilton's principle in 1853. Hamilton's principle became the cornerstone for classical work with different forms of action until Richard Feynman and Julian Schwinger developed quantum action principles. Definitions Expressed in mathematical language, using the calculus of variations, the evolution of a physical system (i.e., how the system actually progresses from one state to another) corresponds to a stationary point (usually, a minimum) of the action. Action has the dimensions of [energy] × [time], and its SI unit is joule-second, which is identical to the unit of angular momentum. Several different definitions of "the action" are in common use in physics. The action is usually an integral over time. However, when the action pertains to fields, it may be integrated over spatial variables as well. In some cases, the action is integrated along the path followed by the physical system. The action is typically represented as an integral over time, taken along the path of the system between the initial time and the final time of the development of the system: where the integrand L is called the Lagrangian. For the action integral to be well-defined, the trajectory has to be bounded in time and space. Action (functional) Most commonly, the term is used for a functional which takes a function of time and (for fields) space as input and returns a scalar. In classical mechanics, the input function is the evolution q(t) of the system between two times t1 and t2, where q represents the generalized coordinates. The action is defined as the integral of the Lagrangian L for an input evolution between the two times: where the endpoints of the evolution are fixed and defined as and . According to Hamilton's principle, the true evolution qtrue(t) is an evolution for which the action is stationary (a minimum, maximum, or a saddle point). This principle results in the equations of motion in Lagrangian mechanics. Abbreviated action (functional) In addition to the action functional, there is another functional called the abbreviated action. In the abbreviated action, the input function is the path followed by the physical system without regard to its parameterization by time. For example, the path of a planetary orbit is an ellipse, and the path of a particle in a uniform gravitational field is a parabola; in both cases, the path does not depend on how fast the particle traverses the path. The abbreviated action (sometime written as ) is defined as the integral of the generalized momenta, for a system Lagrangian along a path in the generalized coordinates : where and are the starting and ending coordinates. According to Maupertuis's principle, the true path of the system is a path for which the abbreviated action is stationary. Hamilton's characteristic function When the total energy E is conserved, the Hamilton–Jacobi equation can be solved with the additive separation of variables: where the time-independent function W(q1, q2, ..., qN) is called Hamilton's characteristic function. The physical significance of this function is understood by taking its total time derivative This can be integrated to give which is just the abbreviated action. Action of a generalized coordinate A variable Jk in the action-angle coordinates, called the "action" of the generalized coordinate qk, is defined by integrating a single generalized momentum around a closed path in phase space, corresponding to rotating or oscillating motion: The corresponding canonical variable conjugate to Jk is its "angle" wk, for reasons described more fully under action-angle coordinates. The integration is only over a single variable qk and, therefore, unlike the integrated dot product in the abbreviated action integral above. The Jk variable equals the change in Sk(qk) as qk is varied around the closed path. For several physical systems of interest, Jk is either a constant or varies very slowly; hence, the variable Jk is often used in perturbation calculations and in determining adiabatic invariants. For example, they are used in the calculation of planetary and satellite orbits. Single relativistic particle When relativistic effects are significant, the action of a point particle of mass m travelling a world line C parametrized by the proper time is If instead, the particle is parametrized by the coordinate time t of the particle and the coordinate time ranges from t1 to t2, then the action becomes where the Lagrangian is Action principles and related ideas Physical laws are frequently expressed as differential equations, which describe how physical quantities such as position and momentum change continuously with time, space or a generalization thereof. Given the initial and boundary conditions for the situation, the "solution" to these empirical equations is one or more functions that describe the behavior of the system and are called equations of motion. Action is a part of an alternative approach to finding such equations of motion. Classical mechanics postulates that the path actually followed by a physical system is that for which the action is minimized, or more generally, is stationary. In other words, the action satisfies a variational principle: the principle of stationary action (see also below). The action is defined by an integral, and the classical equations of motion of a system can be derived by minimizing the value of that integral. The action principle provides deep insights into physics, and is an important concept in modern theoretical physics. Various action principles and related concepts are summarized below. Maupertuis's principle In classical mechanics, Maupertuis's principle (named after Pierre Louis Maupertuis) states that the path followed by a physical system is the one of least length (with a suitable interpretation of path and length). Maupertuis's principle uses the abbreviated action between two generalized points on a path. Hamilton's principal function Hamilton's principle states that the differential equations of motion for any physical system can be re-formulated as an equivalent integral equation. Thus, there are two distinct approaches for formulating dynamical models. Hamilton's principle applies not only to the classical mechanics of a single particle, but also to classical fields such as the electromagnetic and gravitational fields. Hamilton's principle has also been extended to quantum mechanics and quantum field theory—in particular the path integral formulation of quantum mechanics makes use of the concept—where a physical system explores all possible paths, with the phase of the probability amplitude for each path being determined by the action for the path; the final probability amplitude adds all paths using their complex amplitude and phase. Hamilton–Jacobi equation Hamilton's principal function is obtained from the action functional by fixing the initial time and the initial endpoint while allowing the upper time limit and the second endpoint to vary. The Hamilton's principal function satisfies the Hamilton–Jacobi equation, a formulation of classical mechanics. Due to a similarity with the Schrödinger equation, the Hamilton–Jacobi equation provides, arguably, the most direct link with quantum mechanics. Euler–Lagrange equations In Lagrangian mechanics, the requirement that the action integral be stationary under small perturbations is equivalent to a set of differential equations (called the Euler–Lagrange equations) that may be obtained using the calculus of variations. Classical fields The action principle can be extended to obtain the equations of motion for fields, such as the electromagnetic field or gravitational field. Maxwell's equations can be derived as conditions of stationary action. The Einstein equation utilizes the Einstein–Hilbert action as constrained by a variational principle. The trajectory (path in spacetime) of a body in a gravitational field can be found using the action principle. For a free falling body, this trajectory is a geodesic. Conservation laws Implications of symmetries in a physical situation can be found with the action principle, together with the Euler–Lagrange equations, which are derived from the action principle. An example is Noether's theorem, which states that to every continuous symmetry in a physical situation there corresponds a conservation law (and conversely). This deep connection requires that the action principle be assumed. Path integral formulation of quantum field theory In quantum mechanics, the system does not follow a single path whose action is stationary, but the behavior of the system depends on all permitted paths and the value of their action. The action corresponding to the various paths is used to calculate the path integral, which gives the probability amplitudes of the various outcomes. Although equivalent in classical mechanics with Newton's laws, the action principle is better suited for generalizations and plays an important role in modern physics. Indeed, this principle is one of the great generalizations in physical science. It is best understood within quantum mechanics, particularly in Richard Feynman's path integral formulation, where it arises out of destructive interference of quantum amplitudes. Modern extensions The action principle can be generalized still further. For example, the action need not be an integral, because nonlocal actions are possible. The configuration space need not even be a functional space, given certain features such as noncommutative geometry. However, a physical basis for these mathematical extensions remains to be established experimentally.	Physical sciences	Classical mechanics	Physics
312903	https://en.wikipedia.org/wiki/Excavator	Excavator	Excavators are heavy construction equipment primarily consisting of a boom, dipper (or stick), bucket, and cab on a rotating platform known as the "house". The modern excavator's house sits atop an undercarriage with tracks or wheels, being an evolution of the steam shovel (which itself evolved into the power shovel when steam was replaced by diesel and electric power). All excavation-related movement and functions of a hydraulic excavator are accomplished through the use of hydraulic fluid, with hydraulic cylinders and hydraulic motors, which replaced winches, chains, and steel ropes. Another principle change was the direction of the digging action, with modern excavators pulling their buckets toward them like a dragline rather than pushing them away to fill them the way the first powered shovels did. Terminology Excavators are also called diggers, scoopers, mechanical shovels, or 360-degree excavators (sometimes abbreviated simply to "360"). Tracked excavators are sometimes called "trackhoes" by analogy to the backhoe. In the UK, wheeled excavators are sometimes known as "rubber ducks". Usage Excavators are used in many ways: Digging of trenches, holes, foundations Material handling Brush cutting with hydraulic saw, mower, and stump removal attachments Forestry work Forestry mulching Demolition with hydraulic claw, cutter and breaker attachments Mining, especially, but not only open-pit mining River dredging Hydro excavation to access fragile underground infrastructure using high pressure water Driving piles, in conjunction with a pile driver Drilling shafts for footings and rock blasting, by use of an auger or hydraulic drill attachment Snow removal with snowplow and snow blower attachments Aircraft recycling Configurations Modern hydraulic excavators come in a wide variety of sizes. The smaller ones are called mini or compact excavators. For example, Caterpillar's smallest mini-excavator weighs and has 13 hp; their largest model is the largest excavator available (developed and produced by the Orenstein & Koppel, Germany, until the takeover 2011 by Caterpillar, named »RH400«), the CAT 6090, which weighs in excess of , has 4500 hp, and a bucket as large as 52.0 m3. Hydraulic excavators usually couple engine power to (commonly) three hydraulic pumps rather than to mechanical drivetrains. The two main pumps supply oil at high pressure (up to 5000 psi, 345 bar) for the arms, swing motor, track motors and accessories while the third is a lower pressure (≈700 psi, 48 bar) pump for pilot control of the spool valves; this third circuit allows for reduced physical effort when operating the controls. Generally, the 3 pumps used in excavators consist of 2 variable displacement piston pumps and a gear pump. The arrangement of the pumps in the excavator unit changes with different manufacturers using different formats. The three main sections of an excavator are the undercarriage, the house and the arm. The boom, the front part that is attached to the cab itself and holds the arm, is also used. The undercarriage includes tracks, track frame, and final drives, which have a hydraulic motor and gearing providing the drive to the individual tracks. The undercarriage, especially frequently for a mini-excavator, can also have blade similar to that of a bulldozer. The house includes the operator cab, counterweight, engine, fuel and hydraulic oil tanks. The house attaches to the undercarriage by way of a center pin. High-pressure oil is supplied to the tracks' hydraulic motors through a hydraulic swivel at the axis of the pin, allowing the machine to slew 360° unhindered and thus provides the left-and-right movement. The arm provides the up-and-down and closer-and-further (or digging movement) movements. Arms typically consist of a boom, stick and bucket with three joints between them and the house. The boom attaches to the house and provides the up-and-down movement. It can be one of several different configurations: Most common are mono booms; these have no movement apart from straight up and down. Some others have a knuckle boom which can also move left and right in line with the machine. Another option is a hinge at the base of the boom allowing it to hydraulically pivot up to 180° independent to the house; however, this is generally available only to compact excavators. Variable angle booms have additional joint in the middle of the boom to change the curvature of the boom. These are also called triple-articulated booms (TAB) or 3 piece booms. Attached to the end of the boom is the stick (or dipper arm). The stick provides the digging movement needed to pull the bucket through the ground. The stick length is optional depending whether reach (longer stick) or break-out power (shorter stick) is required. Most common is mono stick but there are also, for example, telescopic sticks. The largest form ever of an excavator, the dragline excavator, eliminated the dipper in favor of a line and winch. On the end of the stick is usually a bucket. A wide, large capacity (mud) bucket with a straight cutting edge is used for cleanup and levelling or where the material to be dug is soft, and teeth are not required. A general purpose (GP) bucket is generally smaller, stronger, and has hardened side cutters and teeth used to break through hard ground and rocks. Buckets have numerous shapes and sizes for various applications. There are also many other attachments that are available to be attached to the excavator for boring, ripping, crushing, cutting, lifting, etc. Attachments can be attached with pins similar to other parts of the arm or with some variety of quick coupler. Excavators in Scandinavia often feature a tiltrotator which allows attachments rotate 360 degrees and tilt +/- 45 degrees, in order to increase the flexibility and precision of the excavator. Before the 1990s, all excavators had a long or conventional counterweight that hung off the rear of the machine to provide more digging force and lifting capacity. This became a nuisance when working in confined areas. In 1993 Yanmar launched the world's first Zero Tail Swing excavator, which allows the counterweight to stay inside the width of the tracks as it slews, thus being safer and more user friendly when used in a confined space. This type of machine is now widely used throughout the world. There are two main types of control configuration used in excavators to control the boom and bucket, each distributing the four primary digging functions across two x-y joysticks. This allows a skilled operator to control all four functions simultaneously. The most popular configuration in the US is the SAE controls configuration while in other parts of the world, the ISO control configuration is more common. Some manufacturers such as Takeuchi have switches that allow the operator to select which control configuration to use. Excavator attachments Hydraulic excavators now perform tasks well beyond bucket excavation. With the advent of hydraulic-powered attachments such as a breaker, a cutter, a grapple or an auger,a crusher and screening buckets the excavator is frequently used in many applications other than excavation. Many excavators feature a quick coupler for simplified attachment mounting, increasing the machine's utilization on the jobsite. Excavators are usually employed together with loaders and bulldozers. Most wheeled, compact and some medium-sized (11 to 18-tonne) excavators have a backfill (or dozer) blade. This is a horizontal bulldozer-like blade attached to the undercarriage and is used for leveling and pushing removed material back into a hole. Notable manufacturers Current manufacturers As of July 2021, current excavator manufacturers include:	Technology	Specific-purpose transportation	null
312933	https://en.wikipedia.org/wiki/Tenrec	Tenrec	A tenrec () is a mammal belonging to any species within the afrotherian family Tenrecidae, which is endemic to Madagascar. Tenrecs are a very diverse group; as a result of adaptive radiation and exhibit convergent evolution, some resemble hedgehogs, shrews, opossums, rats, and mice. They occupy aquatic, arboreal, terrestrial, and fossorial environments. Some of these species, including the greater hedgehog tenrec, can be found in the Madagascar dry deciduous forests. However, the speciation rate in this group has been higher in humid forests. All tenrecs are believed to descend from a common ancestor that lived 29–37 million years ago after rafting over from Africa. The split from their closest relatives, African otter shrews, is estimated to have occurred about 47–53 million years ago. Etymology The word tenrec is borrowed, via French, from the Malagasy word (variant of ), which refers to the tailless tenrec (Tenrec ecaudatus); it has been speculated that the Malagasy word is related to . Evolution Tenrecs are believed to have evolved from a single species that colonized Madagascar between 42 and 25 million years ago. The question of how this family reached Madagascar is still unresolved, but the leading hypothesis suggests a small number of individuals may have found themselves on floating vegetation and crossed the Mozambique Channel, which separates Madagascar from southeastern Africa. The Tenrecidae family is one of only four extant terrestrial mammal lineages to have colonized and diversified on Madagascar. Once established on Madagascar, tenrecs diversified to occupy various niches on the island. Many evolved resemblances to familiar but unrelated mammals that are not found on Madagascar. For instance, the two species of hedgehog tenrec possess coats of hardened spines and the ability to roll into a ball when threatened, characteristics similar to those of true hedgehogs. This example, along with others, demonstrates convergent evolution; it has provided evolutionary biologists with opportunities to study adaptation over evolutionary timescales. Characteristics Tenrecs are small mammals of variable body form. The smallest species are the size of shrews, with a body length of around , and weighing just , while the largest, the common or tailless tenrec, is in length, and can weigh over . Although they may resemble shrews, hedgehogs, or opossums, they are not closely related to any of these groups, their closest relatives being the otter shrews, and after that other African insectivorous mammals including golden moles and elephant shrews. The common ancestry of these animals, which are classified together in the clade Afrotheria, was not recognized until the late 1990s. Continuing work on the molecular and morphological diversity of afrotherian mammals has provided ever increasing support for their common ancestry. Tenrecs are among the few terrestrial mammals that echolocate. Unusual among placental mammals, the rectum and urogenital tracts of tenrecs share a common opening, or cloaca which is a feature more commonly seen in birds, reptiles, and amphibians. They have a low body temperature, sufficiently low that they do not require a scrotum to cool their sperm as do most other mammals. All species appear to be at least somewhat omnivorous, with invertebrates forming the largest part of their diets. One species, Microgale mergulus, is semiaquatic (similar to the lifestyle of their closest relatives, the otter shrews). All of the species, semiaquatic or not, appear to have evolved from a single, common ancestor with the otter shrews comprising the next, most-closely related mammalian species. While the fossil record of tenrecs is scarce, at least some specimens from the early Miocene of Kenya show close affinities to living species from Madagascar, such as Geogale aurita. Most species are nocturnal and have poor eyesight. Their other senses are well developed, however and they have especially sensitive whiskers. As with many of their other features, the dental formula of tenrecs varies greatly between species; they can have from 32 to 42 teeth in total. Unusual for mammals, the permanent dentition in tenrecs tends not to completely erupt until well after adult body size has been reached. This is one of several anatomical features shared by elephants, hyraxes, sengis, and golden moles (but apparently not aardvarks), consistent with their descent from a common ancestor. Tenrecs have a gestation period of 50 to 64 days, and give birth to a number of relatively undeveloped young. While the otter shrews have just two young per litter, the tailless tenrec can have as many as 32, and females possess up to 29 teats, more than any other mammal. Some tenrec species are social, living in multigenerational family groups with over a dozen individuals. Interaction with humans In the island nation of Mauritius, and also on the Comoran island of Mayotte, some of the inhabitants eat tenrec meat, although it is difficult to obtain (as it is not sold in shops or markets) and difficult to prepare correctly. The lesser hedgehog tenrec (Echinops telfairi) is one of 16 mammalian species that will have its genome sequenced as part of the Mammalian Genome Project. It is increasingly popular in the pet trade, and in the future may serve as an important model organism in biomedicine, as it is only distantly related to the mice, rats, guinea pigs, and rhesus macaques which comprise the most common research animals. Threats Of the 31 species assessed, 24 (77%) are categorized by the IUCN Red List as Least Concern, 1 species as Data Deficient, 4 species as Vulnerable, and 2 species as Endangered. The conservation status of many tenrec species is of concern due to an increase of threats within the last 50 years. The main threats facing tenrecs include habitat loss due to deforestation, fragmentation and degradation, hunting, incidental capture, and climate change. Slash-and-burn agriculture, as well as commercial logging and mining of metals is negatively affecting tenrec species that inhabit forests. Five of the six threatened Tenrec species are dependent on forest habitats. Conservation As of 2022, conservation of the tenrec population is not being prioritized. Because most tenrecs are dependent on forest habitats, conservation efforts would need to include a focus on reduction in deforestation on Madagascar as well as habitat restoration. Current conservation efforts include that of the Madagascar Ankizy Fund, started by a paleontological team from Stony Brook University to improve access to health care and education facilities for villagers in remote areas of Madagascar. A healthy and educated local human population will, in the long term, benefit the Malagasy fauna, such as tenrecs. Species The three subfamilies, eight genera, and 31 extant species of tenrecs are FAMILY TENRECIDAE Subfamily Geogalinae Genus Geogale Large-eared tenrec (Geogale aurita) Subfamily Oryzorictinae Genus Microgale Short-tailed shrew tenrec (Microgale brevicaudata) Cowan's shrew tenrec (Microgale cowani) Drouhard's shrew tenrec (Microgale drouhardi) Dryad shrew tenrec (Microgale dryas) Pale shrew tenrec (Microgale fotsifotsy) Gracile shrew tenrec (Microgale gracilis) Grandidier's shrew tenrec (Microgale grandidieri) Naked-nosed shrew tenrec (Microgale gymnorhyncha) Jenkins's shrew tenrec (Microgale jenkinsae) Northern shrew tenrec (Microgale jobihely) Lesser long-tailed shrew tenrec (Microgale longicaudata) Microgale macpheei (extinct) Major's long-tailed tenrec (Microgale majori) Web-footed tenrec (Microgale mergulus) Montane shrew tenrec (Microgale monticola) Nasolo's shrew tenrec (Microgale nasoloi) Pygmy shrew tenrec (Microgale parvula) Greater long-tailed shrew tenrec (Microgale principula) Least shrew tenrec (Microgale pusilla) Shrew-toothed shrew tenrec (Microgale soricoides) Taiva shrew tenrec (Microgale taiva) Thomas's shrew tenrec (Microgale thomasi) Genus Nesogale Dobson's shrew tenrec (Nesogale dobsoni) Talazac's shrew tenrec (Nesogale talazaci) Genus Oryzorictes Mole-like rice tenrec (Oryzorictes hova) Four-toed rice tenrec (Oryzorictes tetradactylus) Subfamily Tenrecinae Tribe Setiferini Genus Echinops Lesser hedgehog tenrec (Echinops telfairi) Genus Setifer Greater hedgehog tenrec (Setifer setosus) Tribe Tenrecini Genus Hemicentetes Highland streaked tenrec (Hemicentetes nigriceps) Lowland streaked tenrec (Hemicentetes semispinosus) Genus Tenrec Common tenrec (Tenrec ecaudatus)	Biology and health sciences	Other afrotheres	Animals
312940	https://en.wikipedia.org/wiki/Pinaceae	Pinaceae	The Pinaceae (), or pine family, are conifer trees or shrubs, including many of the well-known conifers of commercial importance such as cedars, firs, hemlocks, piñons, larches, pines and spruces. The family is included in the order Pinales, formerly known as Coniferales. Pinaceae have distinctive cones with woody scales bearing typically two ovules, and are supported as monophyletic by both morphological trait and genetic analysis. They are the largest extant conifer family in species diversity, with between 220 and 250 species (depending on taxonomic opinion) in 11 genera, and the second-largest (after Cupressaceae) in geographical range, found in most of the Northern Hemisphere, with the majority of the species in temperate climates, but ranging from subarctic to tropical. The family often forms the dominant component of boreal, coastal, and montane forests. One species, Pinus merkusii, grows just south of the equator in Southeast Asia. Major centres of diversity are found in the mountains of southwest China, Mexico, central Japan, and California. Description Members of the family Pinaceae are trees (rarely shrubs) growing from tall, mostly evergreen (except the deciduous Larix and Pseudolarix), resinous, monoecious, with subopposite or whorled branches, and spirally arranged, linear (needle-like) leaves. The embryos of Pinaceae have three to 24 cotyledons. The female cones are large and usually woody, long, with numerous spirally arranged scales, and two winged seeds on each scale. The male cones are small, long, and fall soon after pollination; pollen dispersal is by wind. Seed dispersal is mostly by wind, but some species have large seeds with reduced wings, and are dispersed by birds. Analysis of Pinaceae cones reveals how selective pressure has shaped the evolution of variable cone size and function throughout the family. Variation in cone size in the family has likely resulted from the variation of seed dispersal mechanisms available in their environments over time. All Pinaceae with seeds weighing less than 90 milligrams are seemingly adapted for wind dispersal. Pines having seeds larger than 100 mg are more likely to have benefited from adaptations that promote animal dispersal, particularly by birds. Pinaceae that persist in areas where tree squirrels are abundant do not seem to have evolved adaptations for bird dispersal. Boreal conifers have many adaptions for winter. The narrow conical shape of northern conifers, and their downward-drooping limbs help them shed snow, and many of them seasonally alter their biochemistry to make them more resistant to freezing, called "hardening". Classification Classification of the subfamilies and genera of Pinaceae has been subject to debate in the past. Pinaceae ecology, morphology, and history have all been used as the basis for methods of analyses of the family. An 1891 publication divided the family into two subfamilies, using the number and position of resin canals in the primary vascular region of the young taproot as the primary consideration. In a 1910 publication, the family was divided into two tribes based on the occurrence and type of long–short shoot dimorphism. A more recent classification divided the subfamilies and genera based on the consideration of features of ovulate cone anatomy among extant and fossil members of the family. Below is an example of how the morphology has been used to classify Pinaceae. The 11 genera are grouped into four subfamilies, based on the microscopical anatomy and the morphology of the cones, pollen, wood, seeds, and leaves: Subfamily Pinoideae (Pinus): cones are biennial, rarely triennial, with each year's scale-growth distinct, forming an umbo on each scale, the cone scale base is broad, concealing the seeds fully from abaxial (below the phloem vessels) view, the seed is without resin vesicles, the seed wing holds the seed in a pair of claws, leaves have primary stomatal bands adaxial (above the xylem) or equally on both surfaces. Subfamily Piceoideae (Picea): cones are annual, without a distinct umbo, the cone scale base is broad, concealing the seeds fully from abaxial view, seed is without resin vesicles, blackish, the seed wing holds the seed loosely in a cup, leaves have primary stomatal bands adaxial (above the xylem) or equally on both surfaces. Subfamily Laricoideae (Larix, Pseudotsuga, and Cathaya): cones are annual, without a distinct umbo, the cone scale base is broad, concealing the seeds fully from abaxial view, the seed is without resin vesicles, whitish, the seed wing holds the seed tightly in a cup, leaves have primary stomatal bands abaxial only. Subfamily Abietoideae (Abies, Cedrus, Pseudolarix, Keteleeria, Nothotsuga, and Tsuga): cones are annual, without a distinct umbo, the cone scale base is narrow, with the seeds partly visible in abaxial view, the seed has resin vesicles, the seed wing holds the seed tightly in a cup, leaves have primary stomatal bands abaxial only. Phylogeny A revised 2018 phylogeny places Cathaya as sister to the pines rather than in the Laricoidae subfamily with Larix and Pseudotsuga. Multiple molecular studies indicate that in contrast to previous classifications placing it outside the conifers, Gnetophyta may in fact be the sister group to the Pinaceae, with both lineages having diverged during the early-mid Carboniferous. This is known as the "gnepine" hypothesis. Evolutionary history Pinaceae is estimated to have diverged from other conifer groups during the late Carboniferous ~313 million years ago. Various possible stem-group relatives have been reported from as early as the Late Permian (Lopingian) The extinct conifer cone genus Schizolepidopsis likely represent stem-group members of the Pinaceae, the first good records of which are in the Middle-Late Triassic, with abundant records during the Jurassic across Eurasia. The oldest crown group (descendant of the last common ancestor of all living species) member of Pinaceae is the cone Eathiestrobus, known from the Upper Jurassic (lower Kimmeridgian, 157.3-154.7 million years ago) of Scotland, which likely belongs to the pinoid grouping of the family. Pinaceae rapidly radiated during the Early Cretaceous. Members of the modern genera Pinus (pines), Picea (spruce) and Cedrus (cedar) first appear during the Early Cretaceous. The extinct Cretaceous genera Pseudoaraucaria and Obirastrobus appear to be members of Abietoideae, while Pityostrobus appears to be non-monophyletic, containing many disparately related members of Pinaceae. While Pinaceae, and indeed all of its subfamilies, substantially predate the break up of the super-continent Pangea, its distribution was limited to northern Laurasia. During the Cenozoic, Pinaceae had higher rates of species turnover than Southern Hemisphere conifers, thought to be driven by range shifts in response to glacial cycles. Defense mechanisms External stresses on plants have the ability to change the structure and composition of forest ecosystems. Common external stress that Pinaceae experience are herbivore and pathogen attack which often leads to tree death. In order to combat these stresses, trees need to adapt or evolve defenses against these stresses. Pinaceae have evolved myriad mechanical and chemical defenses, or a combination of the two, in order to protect themselves against antagonists. Pinaceae have the ability to up-regulate a combination of constitutive mechanical and chemical strategies to further their defenses. Pinaceae defenses are prevalent in the bark of the trees. This part of the tree contributes a complex defensive boundary against external antagonists. Constitutive and induced defenses are both found in the bark. Constitutive defenses Constitutive defenses are typically the first line of defenses used against antagonists and can include sclerified cells, lignified periderm cells, and secondary compounds such as phenolics and resins. Constitutive defenses are always expressed and offer immediate protection from invaders but could also be defeated by antagonists that have evolved adaptations to these defense mechanisms. One of the common secondary compounds used by Pinaceae are phenolics or polyphenols. These secondary compounds are preserved in vacuoles of polyphenolic parenchyma cells (PP) in the secondary phloem. Induced defenses Induced defense responses need to be activated by certain cues, such as herbivore damage or other biotic signals. A common induced defense mechanism used by Pinaceae is resins. Resins are also one of the primary defenses used against attack. Resins are short term defenses that are composed of a complex combination of volatile mono- (C10) and sesquiterpenes (C15) and nonvolatile diterpene resin acids (C20). They are produced and stored in specialized secretory areas known as resin ducts, resin blisters, or resin cavities. Resins have the ability to wash away, trap, fend off antagonists, and are also involved in wound sealing. They are an effective defense mechanism because they have toxic and inhibitory effects on invaders, such as insects or pathogens. Resins could have developed as an evolutionary defense against bark beetle attacks. One well researched resin present in Pinaceae is oleoresin. Oleoresin had been found to be a valuable part of the conifer defense mechanism against biotic attacks. They are found in secretory tissues in tree stems, roots, and leaves. Oleoresin is also needed in order to classify conifers. Active research: methyl jasmonate The topic of defense mechanisms within family Pinaceae is a very active area of study with numerous studies being conducted. Many of these studies use methyl jasmonate (MJ) as an antagonist. Methyl jasmonate is known to be able to induce defense responses in the stems of multiple Pinaceae species. It has been found that MJ stimulated the activation of PP cells and formation of xylem traumatic resin ducts (TD). These are structures that are involved in the release of phenolics and resins, both forms of defense mechanism.	Biology and health sciences	Pinaceae	Plants
312943	https://en.wikipedia.org/wiki/Blacklight	Blacklight	A blacklight, also called a UV-A light, Wood's lamp, or ultraviolet light, is a lamp that emits long-wave (UV-A) ultraviolet light and very little visible light. One type of lamp has a violet filter material, either on the bulb or in a separate glass filter in the lamp housing, which blocks most visible light and allows through UV, so the lamp has a dim violet glow when operating. Blacklight lamps which have this filter have a lighting industry designation that includes the letters "BLB". This stands for "blacklight blue". A second type of lamp produces ultraviolet but does not have the filter material, so it produces more visible light and has a blue color when operating. These tubes are made for use in "bug zapper" insect traps, and are identified by the industry designation "BL". This stands for "blacklight". Blacklight sources may be specially designed fluorescent lamps, mercury-vapor lamps, light-emitting diodes (LEDs), lasers, or incandescent lamps. In medicine, forensics, and some other scientific fields, such a light source is referred to as a Wood's lamp, named after Robert Williams Wood, who invented the original Wood's glass UV filters. Although many other types of lamp emit ultraviolet light with visible light, blacklights are essential when UV-A light without visible light is needed, particularly in observing fluorescence, the colored glow that many substances emit when exposed to UV. They are employed for decorative and artistic lighting effects, diagnostic and therapeutic uses in medicine, the detection of substances tagged with fluorescent dyes, rock-hunting, scorpion-hunting, the detection of counterfeit money, the curing of plastic resins, attracting insects and the detection of refrigerant leaks affecting refrigerators and air conditioning systems. Strong sources of long-wave ultraviolet light are used in tanning beds. Medical hazard UV-A presents a potential hazard when eyes and skin are exposed, especially to high power sources. According to the World Health Organization, UV-A is responsible for the initial tanning of skin and it contributes to skin ageing and wrinkling. UV-A may also contribute to the progression of skin cancers. Additionally, UV-A can have negative effects on eyes in both the short-term and long-term. Types Fluorescent Fluorescent blacklight tubes are typically made in the same fashion as normal fluorescent tubes except that a phosphor that emits UVA light instead of visible white light is used on the inside of the tube. The type most commonly used for blacklights, designated blacklight blue or "BLB" by the industry, has a dark blue filter coating on the tube, which filters out most visible light, so that fluorescence effects can be observed. These tubes have a dim violet glow when operating. They should not be confused with "blacklight" or "BL" tubes, which have no filter coating, and have a brighter blue color. These are made for use in "bug zapper" insect traps where the emission of visible light does not interfere with the performance of the product. The phosphor typically used for a near 368 to 371 nanometer emission peak is either europium-doped strontium fluoride (:) or europium-doped strontium borate (:) while the phosphor used to produce a peak around 350 to 353 nanometres is lead-doped barium silicate (:). "Blacklight blue" lamps peak at 365 nm. Manufacturers use different numbering systems for blacklight tubes. Philips' is becoming outdated (as of 2010), while the (German) Osram system is becoming dominant outside North America. The following table lists the tubes generating blue, UVA and UVB, in order of decreasing wavelength of the most intense peak. Approximate phosphor compositions, major manufacturer's type numbers and some uses are given as an overview of the types available. "Peak" position is approximated to the nearest 10 nm. "Width" is the measure between points on the shoulders of the peak that represent 50% intensity. Bug zappers Another class of UV fluorescent bulb is designed for use in bug zappers. Insects are attracted to the UV light, which they are able to see, and are then electrocuted by the device. These bulbs use the same UV-A emitting phosphor blend as the filtered blacklight, but since they do not need to suppress visible light output, they do not use a purple filter material in the bulb. Plain glass blocks out less of the visible mercury emission spectrum, making them appear light blue-violet to the naked eye. These lamps are referred to by the designation "blacklight" or "BL" in some North American lighting catalogs. These types are not suitable for applications which require the low visible light output of "BLB" tubes lamps. Incandescent A blacklight may also be formed by simply using a UV filter coating such as Wood's glass on the envelope of a common incandescent bulb. This was the method that was used to create the very first blacklight sources. Although incandescent bulbs are a cheaper alternative to fluorescent tubes, they are exceptionally inefficient at producing UV light since most of the light emitted by the filament is visible light which must be blocked. Due to its black body spectrum, an incandescent light radiates less than 0.1% of its energy as UV light. Incandescent UV bulbs, due to the necessary absorption of the visible light, become very hot during use. This heat is, in fact, encouraged in such bulbs, since a hotter filament increases the proportion of UVA in the black-body radiation emitted. This high running-temperature reduces the life of the lamp from a typical 1,000 hours to around 100 hours. Mercury vapor High-power mercury vapor blacklight lamps are made in power ratings of 100 to 1,000 watts. These do not use phosphors, but rely on the intensified and slightly broadened 350–375 nm spectral line of mercury from high pressure discharge at between , depending upon the specific type. These lamps use envelopes of Wood's glass or similar optical filter coatings to block out all the visible light and also the short wavelength (UVC) lines of mercury at 184.4 and 253.7 nm, which are harmful to the eyes and skin. A few other spectral lines, falling within the pass band of the Wood's glass between 300 and 400 nm, contribute to the output. These lamps are used mainly for theatrical purposes and concert displays. They are more efficient UVA producers per unit of power consumption than fluorescent tubes. LED Ultraviolet light can be generated by some light-emitting diodes, but wavelengths shorter than 380 nm are uncommon, and the emission peaks are broad, so only the very lowest energy UV photons are emitted, within predominant not visible light. Safety Although blacklights produce light in the UV range, their spectrum is mostly confined to the longwave UVA region, that is, UV radiation nearest in wavelength to visible light, with low frequency and therefore relatively low energy. While low, there is still some power of a conventional blacklight in the UVB range. UVA is the safest of the three spectra of UV light, although high exposure to UVA has been linked to the development of skin cancer in humans. The relatively low energy of UVA light does not cause sunburn. It can damage collagen fibers, so may accelerate skin aging and cause wrinkles. It can also degrade vitamin A in the skin. UVA light has been shown to cause DNA damage, but not directly, like UVB and UVC. Due to its longer wavelength, it is absorbed less and reaches deeper into skin layers, where it produces reactive chemical intermediates such as hydroxyl and oxygen radicals, which in turn can damage DNA and result in a risk of melanoma. The weak output of blacklights is not sufficient to cause DNA damage or cellular mutations in the way that direct summer sunlight can, although there are reports that overexposure to the type of UV radiation used for creating artificial suntans on sunbeds can cause DNA damage, photo-aging (damage to the skin from prolonged exposure to sunlight), toughening of the skin, suppression of the immune system, cataract formation and skin cancer. UV-A can have negative effects on eyes in both the short-term and long-term. Uses Ultraviolet radiation is invisible to the human eye, but illuminating certain materials with UV radiation causes the emission of visible light, causing these substances to glow with various colors. This is called fluorescence, and has many practical uses. Blacklights are required to observe fluorescence, since other types of ultraviolet lamps emit visible light which drowns out the dim fluorescent glow. Medical applications A Wood's lamp is a diagnostic tool used in dermatology by which ultraviolet light is shone (at a wavelength of approximately 365 nanometers) onto the skin of the patient; a technician then observes any subsequent fluorescence. For example, porphyrins—associated with some skin diseases—will fluoresce pink. Though the technique for producing a source of ultraviolet light was devised by Robert Williams Wood in 1903 using "Wood's glass", it was in 1925 that the technique was used in dermatology by Margarot and Deveze for the detection of fungal infection of hair. It has many uses, both in distinguishing fluorescent conditions from other conditions and in locating the precise boundaries of the condition. Fungal and bacterial infections It is also helpful in diagnosing: Fungal infections. Some forms of tinea, such as Trichophyton tonsurans, do not fluoresce. Bacterial infections Corynebacterium minutissimum is coral red Pseudomonas is yellow-green Cutibacterium acnes, a bacterium involved in acne causation, exhibits an orange glow under a Wood's lamp. Ethylene glycol poisoning A Wood's lamp may be used to rapidly assess whether an individual is suffering from ethylene glycol poisoning as a consequence of antifreeze ingestion. Manufacturers of ethylene glycol-containing antifreezes commonly add fluorescein, which causes the patient's urine to fluoresce under Wood's lamp. Diagnosis Wood's lamp is useful in diagnosing conditions such as tuberous sclerosis and erythrasma (caused by Corynebacterium minutissimum, see above). Additionally, detection of porphyria cutanea tarda can sometimes be made when urine turns pink upon illumination with Wood's lamp. Wood's lamps have also been used to differentiate hypopigmentation from depigmentation such as with vitiligo. A vitiligo patient's skin will appear yellow-green or blue under the Wood's lamp. Its use in detecting melanoma has been reported. Security and authentication Blacklight is commonly used to authenticate oil paintings, antiques and banknotes. It can also differentiate real currency from counterfeit notes because, in many countries, legal banknotes have fluorescent symbols on them that only show under a blacklight. In addition, the paper used for printing money does not contain any of the brightening agents which cause commercially available papers to fluoresce under blacklight. Both of these features make illegal notes easier to detect and more difficult to successfully counterfeit. The same security features can be applied to identification cards such as passports or driver's licenses. Other security applications include the use of pens containing a fluorescent ink, generally with a soft tip, that can be used to "invisibly" mark items. If the objects that are so marked are subsequently stolen, a blacklight can be used to search for these security markings. At some amusement parks, nightclubs and at other, day-long (or night-long) events, a fluorescent mark is rubber stamped onto the wrist of a guest who can then exercise the option of leaving and being able to return again without paying another admission fee. Biology Fluorescent materials are also very widely used in numerous applications in molecular biology, often as "tags" which bind themselves to a substance of interest (for example, DNA), so allowing their visualization. Thousands of moth and insect collectors all over the world use various types of blacklights to attract moth and insect specimens for photography and collecting. It is one of the preferred light sources for attracting insects and moths at night. They can illuminate animal excreta, such as urine and vomit, that is not always visible to the naked eye. Fault detection Blacklight is used extensively in non-destructive testing. Fluorescing fluids are applied to metal structures and illuminated, allowing easy detection of cracks and other weaknesses. If a leak is suspected in a refrigerator or an air conditioning system, a UV tracer dye can be injected into the system along with the compressor lubricant oil and refrigerant mixture. The system is then run in order to circulate the dye across the piping and components and then the system is examined with a blacklight lamp. Any evidence of fluorescent dye then pinpoints the leaking part which needs replacement. Art and decor Blacklight is used to illuminate pictures painted with fluorescent colors, particularly on black velvet, which intensifies the illusion of self-illumination. The use of such materials, often in the form of tiles viewed in a sensory room under UV light, is common in the United Kingdom for the education of students with profound and multiple learning difficulties. Such fluorescence from certain textile fibers, especially those bearing optical brightener residues, can also be used for recreational effect, as seen, for example, in the opening credits of the James Bond film A View to a Kill. Blacklight puppetry is performed in a blacklight theater. Mineral identification Blacklights are a common tool for rock-hunting and identification of minerals by their fluorescence. The most common minerals and rocks that glow under UV light are fluorite, calcite, aragonite, opal, apatite, chalcedony, corundum (ruby and sapphire), scheelite, selenite, smithsonite, sphalerite, sodalite. The first person to observe fluorescence in minerals was George Stokes in 1852. He noted the ability of fluorite to produce a blue glow when illuminated with ultraviolet light and called this phenomenon “fluorescence” after the mineral fluorite. Lamps used to visualise seams of fluorite and other fluorescent minerals are commonly used in mines but they tend to be on an industrial scale. The lamps need to be short wavelength to be useful for this purpose and of scientific grade. UVP range of hand held UV lamps are ideal for this purpose and are used by Geologists to identify the best sources of fluorite in mines or potential new mines. Some transparent selenite crystals exhibit an “hourglass” pattern under UV light that is not visible in natural light. These crystals are also phosphorescent. Limestone, marble, and travertine can glow because of calcite presence. Granite, syenite, and granitic pegmatite rocks can also glow. Curing resins UV light can be used to harden particular glues, resins and inks by causing a photochemical reaction inside those substances. This process of hardening is called ‘curing’. UV curing is adaptable to printing, coating, decorating, stereolithography, and in the assembly of a variety of products and materials. In comparison to other technologies, curing with UV energy may be considered a low-temperature process, a high-speed process, and is a solventless process, as cure occurs via direct polymerization rather than by evaporation. Originally introduced in the 1960s, this technology has streamlined and increased automation in many industries in the manufacturing sector. A primary advantage of curing with ultraviolet light is the speed at which a material can be processed. Speeding up the curing or drying step in a process can reduce flaws and errors by decreasing time that an ink or coating spends wet. This can increase the quality of a finished item, and potentially allow for greater consistency. Another benefit to decreasing manufacturing time is that less space needs to be devoted to storing items which can not be used until the drying step is finished. Because UV energy has unique interactions with many different materials, UV curing allows for the creation of products with characteristics not achievable via other means. This has led to UV curing becoming fundamental in many fields of manufacturing and technology, where changes in strength, hardness, durability, chemical resistance, and many other properties are required. Cockpit lighting, LSD testing and tanning One of the innovations for night and all-weather flying used by the US, UK, Japan and Germany during World War II was the use of UV interior lighting to illuminate the instrument panel, giving a safer alternative to the radium-painted instrument faces and pointers, and an intensity that could be varied easily and without visible illumination that would give away an aircraft's position. This went so far as to include the printing of charts that were marked in UV-fluorescent inks, and the provision of UV-visible pencils and slide rules such as the E6B. They may also be used to test for LSD, which fluoresces under blacklight while common substitutes such as 25I-NBOMe do not. Strong sources of long-wave ultraviolet light are used in tanning beds.	Physical sciences	Basics	Physics
313009	https://en.wikipedia.org/wiki/Drug%20overdose	Drug overdose	A drug overdose (overdose or OD) is the ingestion or application of a drug or other substance in quantities much greater than are recommended. Typically the term is applied for cases when a risk to health is a potential result. An overdose may result in a toxic state or death. Classification The word "overdose" implies that there is a common safe dosage and usage for the drug; therefore, the term is commonly applied only to drugs, not poisons, even though many poisons as well are harmless at a low enough dosage. Drug overdose is sometimes used as a means to commit suicide, as the result of intentional or unintentional misuse of medication. Intentional misuse leading to overdose can include using prescribed or non-prescribed drugs in excessive quantities in an attempt to produce euphoria. Usage of illicit drugs, in large quantities, or after a period of drug abstinence can also induce overdose. Cocaine and opioid users who inject intravenously can easily overdose accidentally, as the margin between a pleasurable drug sensation and an overdose is small. Unintentional misuse can include errors in dosage caused by failure to read or understand product labels. Accidental overdoses may also be the result of over-prescription, failure to recognize a drug's active ingredient or unwitting ingestion by children. A common unintentional overdose in young children involves multivitamins containing iron. The term 'overdose' is often misused as a descriptor for adverse drug reactions or negative drug interactions due to mixing multiple drugs simultaneously. Signs and symptoms Signs and symptoms of an overdose vary depending on the drug or exposure to toxins. The symptoms can often be divided into differing toxidromes. This can help one determine what class of drug or toxin is causing the difficulties. Symptoms of opioid overdoses include slow breathing, heart rate and pulse. Opioid overdoses can also cause pinpoint pupils, and blue lips and nails due to low levels of oxygen in the blood. A person experiencing an opioid overdose might also have muscle spasms, seizures and decreased consciousness. A person experiencing an opiate overdose usually will not wake up, even if their name is called or they are shaken vigorously. Causes The drugs or toxins that are most frequently involved in overdose and death (grouped by ICD-10): Acute alcohol intoxication (F10) Ethyl alcohol (alcohol) Methanol poisoning Ethylene glycol poisoning Opioid overdose (F11) Among sedative-hypnotics (F13) Barbiturate overdose (T42.3) Benzodiazepine overdose (T42.4) Uncategorized sedative-hypnotics (T42.6) Ethchlorvynol (Placidyl) GHB Glutethimide (Doriden) Methaqualone Ketamine (T41.2) Among stimulants (F14-F15) Cocaine overdose (T40.5) Amphetamine overdose (T43.6) Methamphetamine overdose (T43.6) Among tobacco (F17) Nicotine poisoning (T65.2) Among poly drug use (F19) Drug "cocktails" (speedballs) Medications Aspirin poisoning (T39.0) Paracetamol poisoning (Alone or mixed with oxycodone) Paracetamol toxicity (T39.1) Tricyclic antidepressant overdose (T43.0) Vitamin poisoning Pesticide poisoning (T60) Organophosphate poisoning DDT Inhalants Lithium toxicity Added flavoring Masking undesired taste may impair judgement of the potency, which is a factor in overdosing. For example, lean is usually created as a drinkable mixture, the cough syrup is combined with soft drinks, especially fruit-flavored drinks such as Sprite, Mountain Dew or Fanta, and is typically served in a foam cup. A hard candy, usually a Jolly Rancher, may be added to give the mixture a sweeter flavor. Diagnosis The substance that has been taken may often be determined by asking the person. However, if they will not, or cannot, due to an altered level of consciousness, provide this information, a search of the home or questioning of friends and family may be helpful. Examination for toxidromes, drug testing, or laboratory test may be helpful. Other laboratory test such as glucose, urea and electrolytes, paracetamol levels and salicylate levels are typically done. Negative drug-drug interactions have sometimes been misdiagnosed as an acute drug overdose, occasionally leading to the assumption of suicide. Prevention The distribution of naloxone to injection drug users and other opioid drug users decreases the risk of death from overdose. The Centers for Disease Control and Prevention (CDC) estimates that U.S. programs for drug users and their caregivers prescribing take-home doses of naloxone and training on its utilization are estimated to have prevented 10,000 opioid overdose deaths. Healthcare institution-based naloxone prescription programs have also helped reduce rates of opioid overdose in the U.S. state of North Carolina, and have been replicated in the U.S. military. Nevertheless, scale-up of healthcare-based opioid overdose interventions is limited by providers' insufficient knowledge and negative attitudes towards prescribing take-home naloxone to prevent opioid overdose. Programs training police and fire personnel in opioid overdose response using naloxone have also shown promise in the U.S. Supervised injection sites (also known as overdose prevention centers) have been used to help prevent drug overdoses by offering opioid reversal medications such as naloxone, medical assistance and treatment options. They also provide clean needles to help prevent the spread of diseases like HIV/AIDS and hepatitis. Management Stabilization of the person's airway, breathing, and circulation (ABCs) is the initial treatment of an overdose. Ventilation is considered when there is a low respiratory rate or when blood gases show the person to be hypoxic. Monitoring of the patient should continue before and throughout the treatment process, with particular attention to temperature, pulse, respiratory rate, blood pressure, urine output, electrocardiography (ECG) and O2 saturation. Poison control centers and medical toxicologists are available in many areas to provide guidance in overdoses both to physicians and to the general public. Antidotes Specific antidotes are available for certain overdoses. For example, naloxone is the antidote for opiates such as heroin or morphine. Similarly, benzodiazepine overdoses may be effectively reversed with flumazenil. As a nonspecific antidote, activated charcoal is frequently recommended if available within one hour of the ingestion and the ingestion is significant. Gastric lavage, syrup of ipecac, and whole bowel irrigation are rarely used. Epidemiology and statistics The UN gives a figure of 300,000 deaths per year in the world through drug overdose. 1,015,060 US residents died from drug overdoses from 1968 to 2019. 22 people out of every 100,000 died from drug overdoses in 2019 in the US. From 1999 to Feb 2019 in the United States, more than 770,000 people have died from drug overdoses. In the US around 107,500 people died in the 12-month period ending August 31, 2022, at a rate of 294 deaths per day. 70,630 people died from drug overdoses in 2019. The U.S. drug overdose death rate has gone from 2.5 per 100,000 people in 1968 to 21.5 per 100,000 in 2019. The National Center for Health Statistics reports that 19,250 people died of accidental poisoning in the U.S. in the year 2004 (eight deaths per 100,000 population). In 2008 testimony before a Senate subcommittee, Leonard J. Paulozzi, a medical epidemiologist at the Centers for Disease Control and Prevention said that in 2005 more than 22,000 American people died due to overdoses, and the number is growing rapidly. Paulozzi also testified that all available evidence suggests unintentional overdose deaths are related to the increasing use of prescription drugs, especially opioid painkillers. However, the vast majority of overdoses are also attributable to alcohol. It is very rare for a victim of an overdose to have consumed just one drug. Most overdoses occur when drugs are ingested in combination with alcohol. Drug overdose was the leading cause of injury death in 2013. Among people 25 to 64 years old, drug overdose caused more deaths than motor vehicle traffic crashes. There were 43,982 drug overdose deaths in the United States in 2013. Of these, 22,767 (51.8%) were related to prescription drugs. The 22,767 deaths relating to prescription drug overdose in 2013, 16,235 (71.3%) involved opioid painkillers, and 6,973 (30.6%) involved benzodiazepines. Drug misuse and abuse caused about 2.5 million emergency department (ED) visits in 2011. Of these, more than 1.4 million ED visits were related to prescription drugs. Among those ED visits, 501,207 visits were related to anti-anxiety and insomnia medications, and 420,040 visits were related to opioid analgesics. New CDC data in 2024 demonstrates U.S. drug overdose deaths have significantly declined, marking the potential for the first year with fewer than 100,000 fatalities since 2020. The CDC data shows a nearly 17% drop in reported overdose deaths during the 12 months ending in June, totaling 93,087. This is a notable decrease from the 111,615 deaths recorded in the same period ending in June 2023. While the opioid crisis continues to take a heavy toll, fentanyl remains a major driver, contributing to the majority of these fatalities.	Biology and health sciences	Types	Health
313073	https://en.wikipedia.org/wiki/Acanthuridae	Acanthuridae	Acanthuridae are a family of ray-finned fish which includes surgeonfishes, tangs, and unicornfishes. The family includes about 86 extant species of marine fish living in tropical seas, usually around coral reefs. Many of the species are brightly colored and popular in aquaria. Etymology and taxonomic history The name of the family is derived from the Greek words akantha and oura, which loosely translate to "thorn" and "tail", respectively. This refers to the distinguishing characteristic of the family, the "scalpel" found on the caudal peduncle. In the early 1900s, the family was called Hepatidae. Subfamilies and genera Acanthuridae contains the following extant subfamilies and genera: Subfamily Nasinae Fowler & Bean, 1929 Genus Naso Lacépède, 1801 Subfamily Acanthurinae Bonaparte, 1835 Tribe Acanthurini Bonaparte, 1839 Genus Acanthurus Forsskål 1775 Genus Ctenochaetus Gill, 1884 Tribe Prionurini J. L. B. Smith, 1966 Genus Prionurus Lacépède, 1804 Tribe Zebrasomini Winterbottom, 1993 Genus Paracanthurus Bleeker, 1863 Genus Zebrasoma Swainson, 1839 Evolution and fossil record There are several extinct genera known from fossils dating from the Eocene to Miocene: Eocene genera Proacanthurus Tylerichthys Gazolaichthys Naseus Tauichthys Eorandallius Metacanthurus Oligocene genera Glarithurus Caprovesposus Arambourgthurus ?Eonaso Miocene genera Marosichthys Morphology The distinctive characteristic of the family is that they have scalpel-like modified scales, one or more on either side of the peduncle of the tail. The spines are dangerously sharp and may seriously injure anyone who carelessly handles such a fish. The dorsal, anal, and caudal fins are large, extending for most of the length of the body. The mouths are small and have a single row of teeth adapted to grazing on algae. Surgeonfishes sometimes feed as solitary individuals, but they often travel and feed in schools. Feeding in schools may be a mechanism for overwhelming the highly aggressive defense responses of small territorial damselfishes that vigorously guard small patches of algae on coral reefs. Most species are fairly small, with a maximum length of , but some in the genus Acanthurus, some in the genus Prionurus, and most species in the genus Naso may grow larger; the whitemargin unicornfish (Naso annulatus) is the largest species in the family, reaching a length of up to . These fishes may grow quickly in aquaria, so average growth size and suitability should be checked before adding them to any marine aquarium. A larval acanthurid, known as an acronurus, looks strikingly different from the juvenile and adult forms of the same individual. It is mostly transparent and tends to have a pelagic lifestyle, living in open water for an extended period of time before settling on the ocean bottom near the shore, where it develops into the juvenile and ultimately the adult form. Symbiotic bacteria Acanthurids are the only known hosts of the bacteria of the genus Epulopiscium bacteria. These bacteria affect the digestion of surgeonfishes enabling them to digest the algae in their diet. In the aquarium Tangs are very sensitive to disease in the home aquarium. However, if the tang is fed enough algae and the aquarium is properly maintained disease should not be a problem. It is usually necessary to quarantine the animals for a period before introducing them to the aquarium. Adults range from in length and most grow quickly even in aquaria. When considering a tang for an aquarium it is important to consider the size to which these fish can grow. Larger species such as the popular Pacific blue tang surgeonfish (of Finding Nemo fame), Naso or lipstick tang, lined surgeonfish, Sohal surgeonfish and Atlantic blue tang surgeonfish can grow to and require swimming room and hiding places. Many also suggest adding aggressive tangs to the aquarium last as they are territorial and may fight and possibly kill other fish. Tangs primarily graze on macroalgae from genera such as Caulerpa and Gracilaria, although they have been observed in an aquarium setting to eat meat-based fish foods. A popular technique for aquarists, is to grow macroalgae in a sump or refugium. This technique not only is economically beneficial, but serves to promote enhanced water quality through nitrate absorption. The growth of the algae can then be controlled by feeding it to the tang. Gallery	Biology and health sciences	Acanthomorpha	Animals
313236	https://en.wikipedia.org/wiki/Phytolacca	Phytolacca	Phytolacca is a genus of perennial plants native to North America, South America and East Asia. Some members of the genus are known as pokeweeds or similar names such as pokebush, pokeberry, pokeroot or poke sallet. Other names for species of Phytolacca include inkberry and ombú. The generic name is derived from the Greek word (phyton), meaning "plant," and the Latin word lacca, a red dye. Phytolaccatoxin and phytolaccigenin are present (in the leaves, stems, roots, blossoms, berries etc.) in many species which are poisonous to mammals if not prepared properly. The berries are eaten by birds, which are not affected by the toxin. The small seeds with very hard outer shells remain intact in the digestive system and are eliminated whole. The genus comprises about 25 to 35 species of perennial herbs, shrubs, and trees growing from tall. They have alternate simple leaves, pointed at the end, with entire or crinkled margins; the leaves can be either deciduous or evergreen. The stems are green, pink or red. The flowers are greenish-white to pink, produced in long racemes at the ends of the stems. They develop into globose berries diameter, green at first, ripening dark purple to black. Selected species The following species are accepted by one or more regional floras: Phytolacca acinosa Roxb. – Indian poke. Southern and eastern Asia (syn. P. esculenta Van Houtte, P. latbenia (Moq.) Walter). Black and Judziewicz report it in Dane County, Wisconsin in their 2008 and 2009 books (Wildflowers of Wisconsin and the Great Lakes Region. A Comprehensive Field Guide, first and 2nd editions. ) Phytolacca americana L. – American pokeweed. North America (syn. P. decandra L.) Phytolacca australis Phil. – Western South America Phytolacca bogotensis Kunth – Tropical and subtropical South America (sometimes included in P. icosandra). Phytolacca chilensis Miers – central Chile (possibly synonymous with P. icosandra) Phytolacca dioica L. – Ombú. Subtropical South America. Phytolacca dodecandra L'Hér. – Eastern Africa, Madagascar (syn. P. abyssinica Hoffm.). Phytolacca heterotepala H.Walt. – Mexican pokeweed. Mexico. Phytolacca icosandra L. – Central and South America. Phytolacca japonica Makino – Eastern Asia (syn. P. hunanensis Hand.-Mazz., P. zhejiangensis W.T.Fan). Phytolacca octandra L. – Red inkplant. Subtropical and tropical regions worldwide (sometimes included in P. icosandra). Phytolacca polyandra Batalin – Central and southwest China (syn. P. clavigera W.W.Smith). Phytolacca pruinosa Fenzl – Levantine Pokeweed. Southern Turkey, Cyprus, Lebanon and Syria. Phytolacca rivinoides Kunth & C.D.Bouché – Central and South America. Phytolacca sandwicensis Endl. – Hawaiian Pokeweed. Hawaii. Phytolacca thyrsiflora Fenzl ex J.A.Schmidt – Northern South America. Phytolacca weberbaueri H.Walt. – Yumbi. Peru. Formerly placed here Leea asiatica (L.) Ridsdale (as P. asiatica L.) Terminalia catappa L. (as P. javanica Osbeck) Ecology The ombú (Phytolacca dioica) grows as a tree on the pampas of South America and is one of the few providers of shade on the open grassland. It is a symbol of Uruguay, Argentina and gaucho culture. P. weberbaueri from Peru also grows to tree size. Both species have massively buttressed bases to their trunks, and very soft wood with a high water storage capacity which makes them resistant to grass fires and drought. In the Pacific Northwest of North America, pokeweed is an invasive species. Uses Phytolacca americana (American pokeweed, pokeweed, poke) is used as a folk medicine and as food, although all parts of it must be considered toxic unless, as folk recipes claim, it is "properly prepared." The root is never eaten and cannot be made edible. Poke salad ('poke salat') is considered part of traditional southern U.S. cuisine, where it is cooked three times in three changes of boiling water to remove some of the harmful components. Toxic constituents which have been identified include the alkaloids phytolaccine and phytolaccotoxin, as well as a glycoprotein. Fossil record A Phytolacca-like fossil has been described from the Upper Cretaceous (late Campanian) Cerro del Pueblo Formation, Coahuila, Mexico, it is a permineralized multiple infructescence composed of berries with six locules, each containing a single seed with a curved embryo developed in a curved ovule with pendulous placentation, a berry anatomy that is similar to that of the genus Phytolacca. Though this new plant from Coahuila shares reproductive characters with Phytolacca, the constant number (six) of carpels per fruit and pendulous placentation support the recognition of a new genus, Coahuilacarpon phytolaccoides.	Biology and health sciences	Caryophyllales	Plants
313295	https://en.wikipedia.org/wiki/Cauterization	Cauterization	Cauterization (or cauterisation, or cautery) is a medical practice or technique of burning a part of a body to remove or close off a part of it. It destroys some tissue in an attempt to mitigate bleeding and damage, remove an undesired growth, or minimize other potential medical harm, such as infections when antibiotics are unavailable. The practice was once widespread for treatment of wounds. Its utility before the advent of antibiotics was said to be effective at more than one level: To prevent exsanguination To close amputations Cautery was historically believed to prevent infection, but current research shows that cautery actually increases the risk for infection by causing more tissue damage and providing a more hospitable environment for bacterial growth. Actual cautery refers to the metal device, generally heated to a dull red glow, that a physician applies to produce blisters, to stop bleeding of a blood vessel, and for other similar purposes. The main forms of cauterization used today are electrocautery and chemical cautery—both are, for example, prevalent in cosmetic removal of warts and stopping nosebleeds. Cautery can also mean the branding of a human. Etymology Cauterize is a Middle English word borrowed from the Old French , from Late Latin "to burn or brand with a hot iron", from Ancient Greek (), from (), "burning or branding iron", and (kaiein) "to burn" (of caustic). History Cauterization has been used to stop heavy bleeding since antiquity. The process was described in the Edwin Smith Papyrus and Hippocratic Corpus. It was primarily used to control hemorrhages, especially those resulting from surgery, in ancient Greece. Archigenes recommended cauterization in the event of hemorrhaging wounds, and Leonides of Alexandria described excising breast tumors and cauterizing the resulting wound in order to control bleeding. The Chinese recommends cauterization as a treatment for various ailments, including dog bites. Indigenous peoples of the Americas, ancient Arabs, and Persians also used the technique. Tools used in the ancient cauterization process ranged from heated lances to cauterizing knives. The piece of metal was heated over fire and applied to the wound. Cauterization continued to be used as a common treatment in medieval times. The Babylonian Talmud (redacted in 500 AD), alluding to the practice, states: "... and the effect of the hot iron comes and removes the traces of the stroke." While mainly employed to stop blood loss, it was also used in cases of tooth extraction and as a treatment for mental illness. In the Muslim world, scholars Al-Zahrawi and Avicenna wrote about techniques and instruments used for cauterization. As late as the 20th-century, Bedouins of the Negev in Israel had it as their practice to take the root of the shaggy sparrow-wort (Thymelaea hirsuta), cut the root into splinters lengthwise, burn the splinter in fire, and then apply the red-hot tip of a splinter to the forehead of a person who was ill with ringworm (dermatophytosis). The technique of ligature of the arteries as an alternative to cauterization was later improved and used more effectively by Ambroise Paré. Electrocautery Electrocauterization is the process of destroying tissue (or cutting through soft tissue) using heat conduction from a metal probe heated by electric current. The procedure stops bleeding from small vessels (larger vessels being ligated). Electrocautery applies high frequency alternating current by a unipolar or bipolar method. It can be a continuous waveform to cut tissue, or intermittent to coagulate tissue. The electrically produced heat in this process inherently can do numerous things to the tissue, depending on the waveform and power level, including cauterize, coagulate, cut, and dry (desiccate). Thus electrocautery, electrocoagulation, electrodesiccation, and electrocurettage are closely related and can co-occur in the same procedure when desired. Electrodesiccation and curettage is a common procedure. Unipolar In unipolar cauterization, the physician contacts the tissue with a single small electrode. The circuit's exit point is a large surface area, such as the buttocks, to prevent electrical burns. The amount of heat generated depends on the size of contact area, power setting or frequency of current, duration of application, and waveform. A constant waveform generates more heat than intermittent. The frequency used in cutting the tissue is higher than in coagulation mode. Bipolar Bipolar electrocautery passes the current between two tips of a forceps-like tool. It has the advantage of not disturbing other electrical body rhythms (such as the heart) and also coagulates tissue by pressure. Lateral thermal injury is greater in unipolar than bipolar devices. Electrocauterization is preferable to chemical cauterization, because chemicals can leach into neighbouring tissue and cauterize outside of intended boundaries. Concern has also been raised regarding toxicity of the surgical smoke electrocautery produces. This contains chemicals that, through inhalation, may harm patients or medical staff. Ultrasonic coagulation and ablation systems are also available. Chemical cautery Many chemical reactions can destroy tissue, and some are used routinely in medicine, most commonly to remove small skin lesions such as warts or necrotized tissue, or for hemostasis. Because chemicals can leach into areas not intended for cauterization, laser and electrical methods are preferable where practical. Some cauterizing agents are: Silver nitrate is the active ingredient of the lunar caustic, a stick that traditionally looks like a large match. It is dipped in water and pressed onto the lesion for a few moments. Trichloroacetic acid Cantharidin is an extract of the blister beetle that causes epidermal necrosis and blistering. It is used to treat warts. Nasal cauterization Frequent nosebleeds are most likely caused by an exposed blood vessel in the nose, usually one in Kiesselbach's plexus. Even if the nose is not bleeding at the time, a physician may cauterize it to prevent future bleeding. Cauterization methods include burning the affected area with acid, hot metal, or lasers. Such a procedure is naturally quite painful. Sometimes, a physician uses liquid nitrogen as a less painful alternative, though it is less effective. A physician may apply cocaine in the few countries that allow it for medical use. Cocaine is the only local anesthetic that also produces vasoconstriction, making it ideal for controlling nosebleeds. More modern treatment applies silver nitrate after a local anesthetic. The procedure is generally painless, but after the anesthetic wears off, there may be pain for several days, and the nose may run for up to a week after this treatment. Nasal cauterization can cause empty nose syndrome. Infant circumcision Cauterization has been used for the circumcision of infants in the United States and Canada. The College of Physicians and Surgeons of Manitoba advises against its use in neonatal circumcision. This method of circumcision resulted in several infants having their penises severely burned.	Biology and health sciences	Treatments	Health
313371	https://en.wikipedia.org/wiki/Rip%20current	Rip current	A rip current (or just rip) is a specific type of water current that can occur near beaches where waves break. A rip is a strong, localized, and narrow current of water that moves directly away from the shore by cutting through the lines of breaking waves, like a river flowing out to sea. The force of the current in a rip is strongest and fastest next to the surface of the water. Rip currents can be hazardous to people in the water. Swimmers who are caught in a rip current and who do not understand what is happening, or who may not have the necessary water skills, may panic, or they may exhaust themselves by trying to swim directly against the flow of water. Because of these factors, rip currents are the leading cause of rescues by lifeguards at beaches. In the United States they cause an average of 71 deaths by drowning per year . A rip current is not the same thing as undertow, although some people use that term incorrectly when they are talking about a rip current. Contrary to popular belief, neither rip nor undertow can pull a person down and hold them under the water. A rip simply carries floating objects, including people, out to just beyond the zone of the breaking waves, at which point the current dissipates and releases everything it is carrying. Causes and occurrence A rip current forms because wind and breaking waves push surface water towards the land. This causes a slight rise in the water level along the shore. This excess water will tend to flow back to the open water via the route of least resistance. When there is a local area which is slightly deeper, such as a break in an offshore sand bar or reef, this can allow water to flow offshore more easily, and this will initiate a rip current through that gap. Water that has been pushed up near the beach flows along the shore towards the outgoing rip as "feeder currents". The excess water flows out at a right angle to the beach, in a tight current called the "neck" of the rip. The "neck" is where the flow is most rapid. When the water in the rip current reaches outside of the lines of breaking waves, the flow disperses sideways, loses power, and dissipates in what is known as the "head" of the rip. Rip currents can form by the coasts of oceans, seas, and large lakes, whenever there are waves of sufficient energy. Rip currents often occur on a gradually shelving shore, where breaking waves approach the shore parallel to it, or where underwater topography encourages outflow at one specific area. Baïnes are one of the patterns identified to be producing rip currents. The location of rip currents can be difficult to predict. Some tend to recur always in the same places, but others can appear and disappear suddenly at various locations along the beach. The appearance and disappearance of rip currents is dependent upon the bottom topography and the direction from which the surf and swells are coming. Rip currents occur wherever there is strong longshore variability in wave breaking. This variability may be caused by such features as sandbars, by piers and jetties, and even by crossing wave trains. They are often located in places where there is a gap in a reef, or low area on a sandbar. Rip currents, once they have formed, may deepen the channel through a sandbar. Rip currents are usually quite narrow, but they tend to be more common, wider, and faster, when and where breaking waves are large and powerful. Local underwater topography makes some beaches more likely to have rip currents. A few beaches are notorious in this respect. Although rip tide is a misnomer, in areas of significant tidal range, rip currents may only occur at certain stages of the tide, when the water is shallow enough to cause the waves to break over a sand bar, but deep enough for the broken wave to flow over the bar. In parts of the world with a big difference between high tide and low tide, and where the shoreline shelves gently, the distance between a bar and the shoreline may vary from a few meters to a kilometer or more, depending whether it is high tide or low tide. A fairly common misconception is that rip currents can pull a swimmer down, under the surface of the water. This is not true, and in reality a rip current is strongest close to the surface, as the flow near the bottom is slowed by friction. The surface of a rip current can often appear to be a relatively smooth area of water, without any breaking waves, and this deceptive appearance may cause some beach-goers to believe that it is a suitable place to enter the water. Technical description A more detailed and technical description of rip currents requires understanding the concept of radiation stress. Radiation stress is the force (or momentum flux) that is exerted on the water column by the presence of the wave. When a wave reaches shallow water and shoals, it increases in height prior to breaking. During this increase in height, radiation stress increases, because of the force exerted by the weight of the water that has been pushed upwards. To balance this, the local mean surface level drops. This is known as the setdown. When the wave breaks and starts reducing in height, the radiation stress decreases as the amount of water that is elevated decreases. When this happens, the mean surface level increases — this is known as the setup. In the formation of a rip current, a wave propagates over a sandbar with a gap in it. When this happens, most of the wave breaks on the sandbar, leading to "setup". The part of the wave that propagates over the gap does not break, and the "setdown" continues in that part. Because of this phenomenon, the mean water surface over the rest of the sandbar is higher than that which is over the gap. The result is a strong flow outward through the gap. This strong flow is the rip current. The vorticity and inertia of rip currents have been studied. From a model of the vorticity of a rip current done at Scripps Institute of Oceanography, it was found that as a fast rip current extends away from shallow water, the vorticity of the current increases, and the width of the current decreases. This model acknowledges that friction plays a role and waves are irregular in nature. From data from Sector-Scanning Doppler Sonar at Scripps Institute of Oceanography, it was found that rip currents in La Jolla, California, lasted several minutes, that they reoccurred one to four times per hour, and that they created a wedge with a 45° arch and a radius of 200–400 meters. Visible characteristics Rip currents have a characteristic appearance, and, with some experience, they can be visually identified from the shore before entering the water. This is helpful to lifeguards, swimmers, surfers, boaters, divers and other water users, who may need to avoid a rip, or in some cases make use of the flow. Rip currents often look somewhat like a road or river running straight out to sea. They are easiest to notice and identify when the zone of breaking waves is viewed from a high vantage point. The following are some visual characteristics that can be used to identify a rip: A noticeable break in the pattern of the waves — the water often looks flat at the rip, in contrast to the lines of breaking waves on either side of the rip. A "river" of foam — the surface of the rip sometimes looks foamy, because the current is carrying foam from the surf out to open water. Different color — the rip may differ in color from the surrounding water. It is often more opaque, cloudier, or muddier, and so, depending on the angle of the sun, the rip may show as darker or lighter than the surrounding water. It is sometimes possible to see that foam or floating debris on the surface of the rip is moving out, away from the shore. In contrast, in the surrounding areas of breaking waves, floating objects and foam are being pushed towards the shore. These characteristics are helpful in learning to recognize and understand the nature of rip currents. Learning these signs can enable a person to recognize the presence and position of rips before entering the water, which is an important skill as studies show the majority of people are unable to identify a rip current and therefore unable to identify safe places to swim. In the United States, some beaches have signs created by the National Oceanic and Atmospheric Administration (NOAA) and United States Lifesaving Association, explaining what a rip current is and how to escape one. These signs are titled, "Rip Currents; Break the Grip of the Rip". Two of these signs are shown in the image at the top of this article. Beachgoers can get information from lifeguards, who are always watching for rip currents, and who will move their safety flags so that swimmers can avoid rips. Danger to swimmers Rip currents are a potential source of danger for people in shallow water with breaking waves, whether this is in seas, oceans or large lakes. Rip currents are the proximate cause of 80% of rescues carried out by beach lifeguards. Rip currents typically flow at about . They can be as fast as , which is faster than any human can swim. Most rip currents are fairly narrow, and even the widest rip currents are not very wide. Swimmers can usually exit the rip easily by swimming at a right angle to the flow, parallel to the beach. Swimmers who are unaware of this fact may exhaust themselves trying unsuccessfully to swim directly against the flow. The flow of the current fades out completely at the head of the rip, outside the zone of the breaking waves, so there is a definite limit to how far the swimmer will be taken out to sea by the flow of a rip current. In a rip current, death by drowning occurs when a person has limited water skills and panics, or when a swimmer persists in trying to swim to shore against a strong rip current, and eventually becomes exhausted and drowns. According to the NOAA rip currents caused an average of 71 deaths annually in the United States over the ten years ending in 2022 (with 69 in 2022). A 2013 Australian study found that rips killed more people in Australia than bushfires, floods, cyclones and shark attacks combined. Survival People caught in a rip current may notice that they are moving away from the shore quite rapidly. Often, it is not possible to swim directly back to shore against a rip current, so this is not recommended. Contrary to popular misunderstanding, a rip does not pull a swimmer under the water. It carries the swimmer away from the shore in a narrow band of moving water. A rip current is like a moving treadmill, which the swimmer can get out of quite easily by swimming at a right angle, across the current, i.e. parallel to the shore in either direction. Rip currents are usually not very wide, so getting out of one only takes a few strokes. Once out of the rip current, getting back to shore is not difficult, since waves are breaking, and floating objects, including swimmers, will be pushed by the waves towards the shore. As an alternative, people who are caught in a strong rip can simply relax, either floating or treading water, and allow the current to carry them until it dissipates completely once it is beyond the surf line. Then the person can signal for help, or swim back through the surf, doing so diagonally, away from the rip and towards the shore. It is necessary for coastal swimmers to understand the danger of rip currents, to learn how to recognize them, and how to deal with them. And when possible, it is necessary that people enter the water only in areas where lifeguards are on duty. In a planned trial in a large rip current at Muriwai Beach in New Zealand, an Australian researcher from the School of Biological, Earth and Environmental Sciences, UNSW Sydney found that "just swim to the side" would not work as the rip current was too wide to see its sides, and said that, despite a rescue boat being near, he was unable to relax and not panic. The current took him 300 metres along the beach in a channel feeding the rip current, and then 400 metres offshore at "speeds approaching those of swimming world records". Uses Experienced and knowledgeable water users, including surfers, body boarders, divers, surf lifesavers and kayakers, when they wish to get out beyond the breaking waves, will sometimes use a rip current as a rapid and effortless means of transportation.	Physical sciences	Oceanography	Earth science
313384	https://en.wikipedia.org/wiki/Long%20division	Long division	In arithmetic, long division is a standard division algorithm suitable for dividing multi-digit Hindu-Arabic numerals (positional notation) that is simple enough to perform by hand. It breaks down a division problem into a series of easier steps. As in all division problems, one number, called the dividend, is divided by another, called the divisor, producing a result called the quotient. It enables computations involving arbitrarily large numbers to be performed by following a series of simple steps. The abbreviated form of long division is called short division, which is almost always used instead of long division when the divisor has only one digit. History Related algorithms have existed since the 12th century. Al-Samawal al-Maghribi (1125–1174) performed calculations with decimal numbers that essentially require long division, leading to infinite decimal results, but without formalizing the algorithm. Caldrini (1491) is the earliest printed example of long division, known as the Danda method in medieval Italy, and it became more practical with the introduction of decimal notation for fractions by Pitiscus (1608). The specific algorithm in modern use was introduced by Henry Briggs 1600. Education Inexpensive calculators and computers have become the most common way to solve division problems, eliminating a traditional mathematical exercise and decreasing the educational opportunity to show how to do so by paper and pencil techniques. (Internally, those devices use one of a variety of division algorithms, the faster of which rely on approximations and multiplications to achieve the tasks.) In North America, long division has been especially targeted for de-emphasis or even elimination from the school curriculum by reform mathematics, though it has been traditionally introduced in the 4th, 5th or even 6th grades. Method In English-speaking countries, long division does not use the division slash or division sign symbols but instead constructs a tableau. The divisor is separated from the dividend by a right parenthesis or vertical bar ; the dividend is separated from the quotient by a vinculum (i.e., an overbar). The combination of these two symbols is sometimes known as a long division symbol or division bracket. It developed in the 18th century from an earlier single-line notation separating the dividend from the quotient by a left parenthesis. The process is begun by dividing the left-most digit of the dividend by the divisor. The quotient (rounded down to an integer) becomes the first digit of the result, and the remainder is calculated (this step is notated as a subtraction). This remainder carries forward when the process is repeated on the following digit of the dividend (notated as 'bringing down' the next digit to the remainder). When all digits have been processed and no remainder is left, the process is complete. An example is shown below, representing the division of 500 by 4 (with a result of 125). 125 (Explanations) 4)500 4 ( 4 × 1 = 4) 10 ( 5 - 4 = 1) 8 ( 4 × 2 = 8) 20 (10 - 8 = 2) 20 ( 4 × 5 = 20) 0 (20 - 20 = 0) A more detailed breakdown of the steps goes as follows: Find the shortest sequence of digits starting from the left end of the dividend, 500, that the divisor 4 goes into at least once. In this case, this is simply the first digit, 5. The largest number that the divisor 4 can be multiplied by without exceeding 5 is 1, so the digit 1 is put above the 5 to start constructing the quotient. Next, the 1 is multiplied by the divisor 4, to obtain the largest whole number that is a multiple of the divisor 4 without exceeding the 5 (4 in this case). This 4 is then placed under and subtracted from the 5 to get the remainder, 1, which is placed under the 4 under the 5. Afterwards, the first as-yet unused digit in the dividend, in this case the first digit 0 after the 5, is copied directly underneath itself and next to the remainder 1, to form the number 10. At this point the process is repeated enough times to reach a stopping point: The largest number by which the divisor 4 can be multiplied without exceeding 10 is 2, so 2 is written above as the second leftmost quotient digit. This 2 is then multiplied by the divisor 4 to get 8, which is the largest multiple of 4 that does not exceed 10; so 8 is written below 10, and the subtraction 10 minus 8 is performed to get the remainder 2, which is placed below the 8. The next digit of the dividend (the last 0 in 500) is copied directly below itself and next to the remainder 2 to form 20. Then the largest number by which the divisor 4 can be multiplied without exceeding 20, which is 5, is placed above as the third leftmost quotient digit. This 5 is multiplied by the divisor 4 to get 20, which is written below and subtracted from the existing 20 to yield the remainder 0, which is then written below the second 20. At this point, since there are no more digits to bring down from the dividend and the last subtraction result was 0, we can be assured that the process finished. If the last remainder when we ran out of dividend digits had been something other than 0, there would have been two possible courses of action: We could just stop there and say that the dividend divided by the divisor is the quotient written at the top with the remainder written at the bottom, and write the answer as the quotient followed by a fraction that is the remainder divided by the divisor. We could extend the dividend by writing it as, say, 500.000... and continue the process (using a decimal point in the quotient directly above the decimal point in the dividend), in order to get a decimal answer, as in the following example. 31.75 4)127.00 12 (12 ÷ 4 = 3) 07 (0 remainder, bring down next figure) 4 (7 ÷ 4 = 1 r 3) 3.0 (bring down 0 and the decimal point) 2.8 (7 × 4 = 28, 30 ÷ 4 = 7 r 2) 20 (an additional zero is brought down) 20 (5 × 4 = 20) 0 In this example, the decimal part of the result is calculated by continuing the process beyond the units digit, "bringing down" zeros as being the decimal part of the dividend. This example also illustrates that, at the beginning of the process, a step that produces a zero can be omitted. Since the first digit 1 is less than the divisor 4, the first step is instead performed on the first two digits 12. Similarly, if the divisor were 13, one would perform the first step on 127 rather than 12 or 1. Basic procedure for long division of Find the location of all decimal points in the dividend and divisor . If necessary, simplify the long division problem by moving the decimals of the divisor and dividend by the same number of decimal places, to the right (or to the left), so that the decimal of the divisor is to the right of the last digit. When doing long division, keep the numbers lined up straight from top to bottom under the tableau. After each step, be sure the remainder for that step is less than the divisor. If it is not, there are three possible problems: the multiplication is wrong, the subtraction is wrong, or a greater quotient is needed. In the end, the remainder, , is added to the growing quotient as a fraction, . Invariant property and correctness The basic presentation of the steps of the process (above) focus on the what steps are to be performed, rather than the properties of those steps that ensure the result will be correct (specifically, that q × m + r = n, where q is the final quotient and r the final remainder). A slight variation of presentation requires more writing, and requires that we change, rather than just update, digits of the quotient, but can shed more light on why these steps actually produce the right answer by allowing evaluation of q × m + r at intermediate points in the process. This illustrates the key property used in the derivation of the algorithm (below). Specifically, we amend the above basic procedure so that we fill the space after the digits of the quotient under construction with 0's, to at least the 1's place, and include those 0's in the numbers we write below the division bracket. This lets us maintain an invariant relation at every step: q × m + r = n, where q is the partially-constructed quotient (above the division bracket) and r the partially-constructed remainder (bottom number below the division bracket). Note that, initially q=0 and r=n, so this property holds initially; the process reduces r and increases q with each step, eventually stopping when r<m if we seek the answer in quotient + integer remainder form. Revisiting the 500 ÷ 4 example above, we find 125 (q, changes from 000 to 100 to 120 to 125 as per notes below) 4)500 400 ( 4 × 100 = 400) 100 (500 - 400 = 100; now q=100, r=100; note q×4+r = 500.) 80 ( 4 × 20 = 80) 20 (100 - 80 = 20; now q=120, r= 20; note q×4+r = 500.) 20 ( 4 × 5 = 20) 0 ( 20 - 20 = 0; now q=125, r= 0; note q×4+r = 500.) Example with multi-digit divisor A divisor of any number of digits can be used. In this example, 1260257 is to be divided by 37. First the problem is set up as follows: 37)1260257 Digits of the number 1260257 are taken until a number greater than or equal to 37 occurs. So 1 and 12 are less than 37, but 126 is greater. Next, the greatest multiple of 37 less than or equal to 126 is computed. So 3 × 37 = 111 < 126, but 4 × 37 > 126. The multiple 111 is written underneath the 126 and the 3 is written on the top where the solution will appear: 3 37)1260257 111 Note carefully which place-value column these digits are written into. The 3 in the quotient goes in the same column (ten-thousands place) as the 6 in the dividend 1260257, which is the same column as the last digit of 111. The 111 is then subtracted from the line above, ignoring all digits to the right: 3 37)1260257 111 15 Now the digit from the next smaller place value of the dividend is copied down and appended to the result 15: 3 37)1260257 111 150 The process repeats: the greatest multiple of 37 less than or equal to 150 is subtracted. This is 148 = 4 × 37, so a 4 is added to the top as the next quotient digit. Then the result of the subtraction is extended by another digit taken from the dividend: 34 37)1260257 111 150 148 22 The greatest multiple of 37 less than or equal to 22 is 0 × 37 = 0. Subtracting 0 from 22 gives 22, we often don't write the subtraction step. Instead, we simply take another digit from the dividend: 340 37)1260257 111 150 148 225 The process is repeated until 37 divides the last line exactly: 34061 37)1260257 111 150 148 225 222 37 Mixed mode long division For non-decimal currencies (such as the British £sd system before 1971) and measures (such as avoirdupois) mixed mode division must be used. Consider dividing 50 miles 600 yards into 37 pieces: mi - yd - ft - in 1 - 634 1 9 r. 15" 37) 50 - 600 - 0 - 0 37 22880 66 348 13 23480 66 348 1760 222 37 333 22880 128 29 15 ===== 111 348 == 170 === 148 22 66 == Each of the four columns is worked in turn. Starting with the miles: 50/37 = 1 remainder 13. No further division is possible, so perform a long multiplication by 1,760 to convert miles to yards, the result is 22,880 yards. Carry this to the top of the yards column and add it to the 600 yards in the dividend giving 23,480. Long division of 23,480 / 37 now proceeds as normal yielding 634 with remainder 22. The remainder is multiplied by 3 to get feet and carried up to the feet column. Long division of the feet gives 1 remainder 29 which is then multiplied by twelve to get 348 inches. Long division continues with the final remainder of 15 inches being shown on the result line. Interpretation of decimal results When the quotient is not an integer and the division process is extended beyond the decimal point, one of two things can happen: The process can terminate, which means that a remainder of 0 is reached; or A remainder could be reached that is identical to a previous remainder that occurred after the decimal points were written. In the latter case, continuing the process would be pointless, because from that point onward the same sequence of digits would appear in the quotient over and over. So a bar is drawn over the repeating sequence to indicate that it repeats forever (i.e., every rational number is either a terminating or repeating decimal). Notation in non-English-speaking countries China, Japan, Korea use the same notation as English-speaking nations including India. Elsewhere, the same general principles are used, but the figures are often arranged differently. Latin America In Latin America (except Argentina, Bolivia, Mexico, Colombia, Paraguay, Venezuela, Uruguay and Brazil), the calculation is almost exactly the same, but is written down differently as shown below with the same two examples used above. Usually the quotient is written under a bar drawn under the divisor. A long vertical line is sometimes drawn to the right of the calculations. 500 ÷ 4 = 125 (Explanations) 4 ( 4 × 1 = 4) 10 ( 5 - 4 = 1) 8 ( 4 × 2 = 8) 20 (10 - 8 = 2) 20 ( 4 × 5 = 20) 0 (20 - 20 = 0) and 127 ÷ 4 = 31.75 124 30 (bring down 0; decimal to quotient) 28 (7 × 4 = 28) 20 (an additional zero is added) 20 (5 × 4 = 20) 0 In Mexico, the English-speaking world notation is used, except that only the result of the subtraction is annotated and the calculation is done mentally, as shown below: 125 (Explanations) 4)500 10 ( 5 - 4 = 1) 20 (10 - 8 = 2) 0 (20 - 20 = 0) In Bolivia, Brazil, Paraguay, Venezuela, French-speaking Canada, Colombia, and Peru, the European notation (see below) is used, except that the quotient is not separated by a vertical line, as shown below: 127\|4 −124 31,75 30 −28 20 −20 0 Same procedure applies in Mexico, Uruguay and Argentina, only the result of the subtraction is annotated and the calculation is done mentally. Eurasia In Spain, Italy, France, Portugal, Lithuania, Romania, Turkey, Greece, Belgium, Belarus, Ukraine, and Russia, the divisor is to the right of the dividend, and separated by a vertical bar. The division also occurs in the column, but the quotient (result) is written below the divider, and separated by the horizontal line. The same method is used in Iran, Vietnam, and Mongolia. 127\|4 −124\|31,75 30 −28 20 −20 0 In Cyprus, as well as in France, a long vertical bar separates the dividend and subsequent subtractions from the quotient and divisor, as in the example below of 6359 divided by 17, which is 374 with a remainder of 1. 6359\|17 −51 \|374 125 \| −119 \| 69\| −68\| 1\| Decimal numbers are not divided directly, the dividend and divisor are multiplied by a power of ten so that the division involves two whole numbers. Therefore, if one were dividing 12,7 by 0,4 (commas being used instead of decimal points), the dividend and divisor would first be changed to 127 and 4, and then the division would proceed as above. In Austria, Germany and Switzerland, the notational form of a normal equation is used. <dividend> : <divisor> = <quotient>, with the colon ":" denoting a binary infix symbol for the division operator (analogous to "/" or "÷"). In these regions the decimal separator is written as a comma. (cf. first section of Latin American countries above, where it's done virtually the same way): 127 : 4 = 31,75 −12 07 −4 30 −28 20 −20 0 The same notation is adopted in Denmark, Norway, Bulgaria, North Macedonia, Poland, Croatia, Slovenia, Hungary, Czech Republic, Slovakia, Vietnam and in Serbia. In the Netherlands, the following notation is used: 12 / 135 \ 11,25 12 15 12 30 24 60 60 0 In Finland, the Italian method detailed above was replaced by the Anglo-American one in the 1970s. In the early 2000s, however, some textbooks have adopted the German method as it retains the order between the divisor and the dividend. Algorithm for arbitrary base Every natural number can be uniquely represented in an arbitrary number base as a sequence of digits where for all , where is the number of digits in . The value of in terms of its digits and the base is Let be the dividend and be the divisor, where is the number of digits in . If , then quotient and remainder . Otherwise, we iterate from , before stopping. For each iteration , let be the quotient extracted so far, be the intermediate dividend, be the intermediate remainder, be the next digit of the original dividend, and be the next digit of the quotient. By definition of digits in base , . By definition of remainder, . All values are natural numbers. We initiate the first digits of . With every iteration, the three equations are true: There only exists one such such that . The final quotient is and the final remainder is Examples In base 10, using the example above with and , the initial values and . Thus, and . In base 16, with and , the initial values are and . Thus, and . If one doesn't have the addition, subtraction, or multiplication tables for base memorised, then this algorithm still works if the numbers are converted to decimal and at the end are converted back to base . For example, with the above example, and with . The initial values are and . Thus, and . This algorithm can be done using the same kind of pencil-and-paper notations as shown in above sections. d8f45 r. 5 12 ) f412df ea a1 90 112 10e 4d 48 5f 5a 5 Rational quotients If the quotient is not constrained to be an integer, then the algorithm does not terminate for . Instead, if then by definition. If the remainder is equal to zero at any iteration, then the quotient is a -adic fraction, and is represented as a finite decimal expansion in base positional notation. Otherwise, it is still a rational number but not a -adic rational, and is instead represented as an infinite repeating decimal expansion in base positional notation. Binary division Performance On each iteration, the most time-consuming task is to select . We know that there are possible values, so we can find using comparisons. Each comparison will require evaluating . Let be the number of digits in the dividend and be the number of digits in the divisor . The number of digits in . The multiplication of is therefore , and likewise the subtraction of . Thus it takes to select . The remainder of the algorithm are addition and the digit-shifting of and to the left one digit, and so takes time and in base , so each iteration takes , or just . For all digits, the algorithm takes time , or in base . Generalizations Rational numbers Long division of integers can easily be extended to include non-integer dividends, as long as they are rational. This is because every rational number has a recurring decimal expansion. The procedure can also be extended to include divisors which have a finite or terminating decimal expansion (i.e. decimal fractions). In this case the procedure involves multiplying the divisor and dividend by the appropriate power of ten so that the new divisor is an integer – taking advantage of the fact that a ÷ b = (ca) ÷ (cb) – and then proceeding as above. Polynomials A generalised version of this method called polynomial long division is also used for dividing polynomials (sometimes using a shorthand version called synthetic division).	Mathematics	Basics	null
313398	https://en.wikipedia.org/wiki/Wingtip%20device	Wingtip device	Wingtip devices are intended to improve the efficiency of fixed-wing aircraft by reducing drag. Although there are several types of wing tip devices which function in different manners, their intended effect is always to reduce an aircraft's drag. Wingtip devices can also improve aircraft handling characteristics and enhance safety for following aircraft. Such devices increase the effective aspect ratio of a wing without greatly increasing the wingspan. Extending the span would lower lift-induced drag, but would increase parasitic drag and would require boosting the strength and weight of the wing. At some point, there is no net benefit from further increased span. There may also be operational considerations that limit the allowable wingspan (e.g., available width at airport gates). Wingtip devices help prevent the flow around the wingtip of higher pressure air under the wing flowing to the lower pressure surface on top at the wingtip, which results in a vortex caused by the forward motion of the aircraft. Winglets also reduce the lift-induced drag caused by wingtip vortices and improve lift-to-drag ratio. This increases fuel efficiency in powered aircraft and increases cross-country speed in gliders, in both cases increasing range. U.S. Air Force studies indicate that a given improvement in fuel efficiency correlates directly with the causal increase in the aircraft's lift-to-drag ratio. Early history Wing end-plates The initial concept dates back to 1897, when English engineer Frederick W. Lanchester patented wing end-plates as a method for controlling wingtip vortices. In the United States, Scottish-born engineer William E. Somerville patented the first functional winglets in 1910. Somerville installed the devices on his early biplane and monoplane designs. Vincent Burnelli received US Patent no: 1,774,474 for his "Airfoil Control Means" on August 26, 1930. Simple flat end-plates did not cause a reduction in drag, because the increase in profile drag was greater than the decrease in induced drag. Hoerner wing tips Following the end of World War II, Dr. Sighard F. Hoerner was a pioneer researcher in the field, having written a technical paper published in 1952 that called for drooped wingtips whose pointed rear tips focused the resulting wingtip vortex away from the upper wing surface. Drooped wingtips are often called "Hoerner tips" in his honor. Gliders and light aircraft have made use of Hoerner tips for many years. The earliest-known implementation of a Hoerner-style downward-angled "wingtip device" on a jet aircraft was during World War II. This was the so-called "Lippisch-Ohren" (Lippisch-ears), allegedly attributed to the Messerschmitt Me 163's designer Alexander Lippisch, and first added to the M3 and M4 third and fourth prototypes of the Heinkel He 162A Spatz jet light fighter for evaluation. This addition was done in order to counteract the dutch roll characteristic present in the original He 162 design, related to its wings having a marked dihedral angle. This became a standard feature of the approximately 320 completed He 162A jet fighters built, with hundreds more He 162A airframes going unfinished by V-E Day. Winglet The term "winglet" was previously used to describe an additional lifting surface on an aircraft, like a short section between wheels on fixed undercarriage. Richard Whitcomb's research in the 1970s at NASA first used winglet with its modern meaning referring to near-vertical extension of the wing tips. The upward angle (or cant) of the winglet, its inward or outward angle (or toe), as well as its size and shape are critical for correct performance and are unique in each application. The wingtip vortex, which rotates around from below the wing, strikes the cambered surface of the winglet, generating a force that angles inward and slightly forward, analogous to a sailboat sailing close hauled. The winglet converts some of the otherwise-wasted energy in the wingtip vortex to an apparent thrust. This small contribution can be worthwhile over the aircraft's lifetime, provided the benefit offsets the cost of installing and maintaining the winglets. Another potential benefit of winglets is that they reduce the intensity of wake vortices. Those trail behind the plane and pose a hazard to other aircraft. Minimum spacing requirements between aircraft operations at airports are largely dictated by these factors. Aircraft are classified by weight (e.g. "Light", "Heavy", etc.) because the vortex strength grows with the aircraft lift coefficient, and thus, the associated turbulence is greatest at low speed and high weight, which produced a high angle of attack. Winglets and wingtip fences also increase efficiency by reducing vortex interference with laminar airflow near the tips of the wing, by 'moving' the confluence of low-pressure (over wing) and high-pressure (under wing) air away from the surface of the wing. Wingtip vortices create turbulence, originating at the leading edge of the wingtip and propagating backwards and inboard. This turbulence 'delaminates' the airflow over a small triangular section of the outboard wing, which destroys lift in that area. The fence/winglet drives the area where the vortex forms upward away from the wing surface, since the center of the resulting vortex is now at the tip of the winglet. The fuel economy improvement from winglets increases with the mission length. Blended winglets allow a steeper angle of attack reducing takeoff distance. Early development Richard T. Whitcomb, an engineer at NASA's Langley Research Center, further developed Hoerner's concept in response to the sharp increase in the cost of fuel after the 1973 oil crisis. With careful aeronautical design he showed that, for a given bending moment, a near-vertical winglet offers a greater drag reduction compared to a horizontal span extension. Whitcomb's designs were flight-tested in 1979–80 by a joint NASA/Air Force team, using a KC-135 Stratotanker based at the Dryden Flight Research Center. A Lockheed L-1011 and McDonnell Douglas DC-10 were also used for testing, and the latter design was directly implemented by McDonnell Douglas on the derivative MD-11, which was rolled out in 1990. In May 1983, a high school student at Bowie High School in Maryland won a grand prize at the 34th International Science and Engineering Fair in Albuquerque, New Mexico for the result of his research on wingtip devices to reduce drag. The same month, he filed a U.S. patent for "wingtip airfoils", published in 1986. Applications NASA NASA's most notable application of wingtip devices is on the Boeing 747 Shuttle Carrier Aircraft. Located on the 747's horizontal stabilizers, the devices increase the tailplane's effectiveness under the weight of the Space Shuttle orbiter, though these were more for directional stability than for drag reduction. Business aircraft Learjet exhibited the prototype Learjet 28 at the 1977 National Business Aviation Association convention. It employed the first winglets ever used on a production aircraft, either civilian or military. Learjet developed the winglet design without NASA assistance. Although the Model 28 was intended to be a prototype experimental aircraft, performance was such that it resulted in a production commitment from Learjet. Flight tests showed that the winglets increased range by about 6.5 percent and improved directional stability. Learjet's application of winglets to production aircraft continued with newer models including the Learjet 55, 31, 60, 45, and Learjet 40. Gulfstream Aerospace explored winglets in the late 1970s and incorporated winglets in the Gulfstream III, Gulfstream IV and Gulfstream V. The Gulfstream V range of allows nonstop routes such as New York–Tokyo, it holds over 70 world and national flight records. The Rutan combined winglets-vertical stabilizer appeared on his Beechcraft Starship business aircraft design that first flew in 1986. Winglets are also applied to other business aircraft, reducing take-off distance to operate from smaller airports, and allowing higher cruise altitudes. Along winglets on new designs, aftermarket vendors developed retrofits. Winglet Technology, LLC of Wichita, Kansas should have tested its elliptical winglets designed to increase payload-range on hot and high departures to retrofit the Citation X. Experimental Conventional winglets were fitted to Rutan's Rutan Voyager, the first aircraft to circumnavigate the world without refueling in 1986. The aircraft's wingtips were damaged, however, when they dragged along the runway during takeoff, removing about from each wingtip, so the flight was made without benefit of winglets. Airliner fuel efficiency The average commercial jet sees a 4-6 percent increase in fuel efficiency and as much as a 6% decrease in in-flight noise from the use of winglets. Actual fuel savings and the related carbon output can vary significantly by plane, route and flight conditions. Wingtip fence A wingtip fence refers to the winglets including surfaces extending both above and below the wingtip, as described in Whitcomb's early research. Both surfaces are shorter than or equivalent to a winglet possessing similar aerodynamic benefits. The Airbus A310-300 was the first airliner with wingtip fences in 1985. Other Airbus models followed with the A300-600, the A320ceo, and the A380. Other Airbus models including the Airbus A320 Enhanced, A320neo, A350 and A330neo have blended winglets rather than wingtip fences. The Antonov An-158 uses wingtip fences. Canted winglets Boeing announced a new version of the 747, the 747-400, in 1985, with an extended range and capacity, using a combination of winglets and increased span to carry the additional load. The winglets increased the 747-400's range by 3.5% over the 747-300, which is otherwise aerodynamically identical but has no winglets. The 747-400D variant lacks the wingtip extensions and winglets included on other 747-400s since winglets would provide minimal benefits on short-haul routes while adding extra weight and cost, although the -400D may be converted to the long-range version if needed. Winglets are preferred for Boeing derivative designs based on existing platforms, because they allow maximum re-use of existing components. Newer designs are favoring increased span, other wingtip devices or a combination of both, whenever possible. The Ilyushin Il-96 was the first Russian and modern jet to feature winglets in 1988. The Bombardier CRJ-100/200 was the first regional airliner to feature winglets in 1992. The A340/A330 followed with canted winglets in 1993/1994. The Tupolev Tu-204 was the first narrowbody aircraft to feature winglets in 1994. The Airbus A220 (née CSeries), from 2016, has canted winglets. Blended winglets A blended winglet is attached to the wing with a smooth curve instead of a sharp angle and is intended to reduce interference drag at the wing/winglet junction. A sharp interior angle in this region can interact with the boundary layer flow causing a drag inducing vortex, negating some of the benefit of the winglet. Seattle-based Aviation Partners develops blended winglets as retrofits for the Gulfstream II, Hawker 800 and the Falcon 2000. On February 18, 2000, blended winglets were announced as an option for the Boeing 737-800; the first shipset was installed on 14 February 2001 and entered revenue service with Hapag-Lloyd Flug on 8 May 2001. The Aviation Partners/Boeing extensions decrease fuel consumption by 4% for long-range flights and increase range by for the 737-800 or the derivative Boeing Business Jet as standard. Also offered for the 737 Classic, many operators have retrofitted their fleets with these for the fuel savings. Aviation Partners Boeing also offers blended winglets for the 757 and 767-300ER. In 2006 Airbus tested two candidate blended winglets, designed by Winglet Technology and Airbus for the Airbus A320 family. In 2009 Airbus launched its "Sharklet" blended winglet, designed to enhance the payload-range of its A320 family and reduce fuel burn by up to 4% over longer sectors. This corresponds to an annual CO2 reduction of 700 tonnes per aircraft. The A320s fitted with Sharklets were delivered beginning in 2012. They are used on the A320neo, the A330neo and the A350. They are also offered as a retrofit option. Raked wingtip Raked wingtips, where the tip has a greater wing sweep than the rest of the wing, are featured on some Boeing Commercial Airplanes to improve fuel efficiency, takeoff and climb performance. Like winglets, they increase the effective wing aspect ratio and diminish wingtip vortices, decreasing lift-induced drag. In testing by Boeing and NASA, they reduce drag by as much as 5.5%, compared to 3.5% to 4.5% for conventional winglets. While an increase in span would be more effective than a same-length winglet, its bending moment is greater. A winglet gives the performance gain of a span increase but has the bending force of a span increase. Raked wingtips offer several weight-reduction advantages relative to simply extending the conventional main wingspan. At high load-factor structural design conditions, the smaller chords of the wingtip are subjected to less load, and they result in less induced loading on the outboard main wing. Additionally, the leading-edge sweep results in the center of pressure being located farther aft than for simple extensions of the span of conventional main wings. At high load factors, this relative aft location of the center of pressure causes the raked wingtip to be twisted more leading-edge down, reducing the bending moment on the inboard wing. However, the relative aft-movement of the center of pressure accentuates flutter. Raked wingtips are installed on the Boeing 767-400ER (first flight on October 9, 1999), all generations of Boeing 777 (June 12, 1994) including the upcoming 777X, the 737-derived Boeing P-8 Poseidon (25 April 2009), all variants of the Boeing 787 (December 15, 2009) (the cancelled Boeing 787-3 would have had a wingspan to fit in ICAO Aerodrome Reference Code D, as its wingspan was decreased by using blended winglets instead of raked wingtips ), and the Boeing 747-8 (February 8, 2010). The Embraer E-jet E2 and C-390 Millennium wings also have raked wingtips. Split-tip The McDonnell Douglas MD-11 was the first aircraft with split-tip winglets in 1990. For the 737 Next Generation, third-party vendor Aviation Partners has introduced a similar design to the 737 MAX wingtip device known as the split scimitar winglet, with United Airlines as the launch customer. The Boeing 737 MAX uses a new type of wingtip device. Resembling a three-way hybrid of a winglet, wingtip fence, and raked wingtip, Boeing claims that this new design should deliver an additional 1.5% improvement in fuel economy over the 10-12% improvement already expected from the 737 MAX. Gliders In 1987, mechanical engineer Peter Masak called on aerodynamicist Mark D. Maughmer, an associate professor of aerospace engineering at the Pennsylvania State University, about designing winglets to improve performance on his wingspan racing sailplane. Others had attempted to apply Whitcomb's winglets to gliders before, and they did improve climb performance, but this did not offset the parasitic drag penalty in high-speed cruise. Masak was convinced it was possible to overcome this hurdle. By trial and error, they ultimately developed successful winglet designs for gliding competitions, using a new PSU–90–125 airfoil, designed by Maughmer specifically for the winglet application. At the 1991 World Gliding Championships in Uvalde, Texas, the trophy for the highest speed went to a winglet-equipped 15-meter class limited wingspan glider, exceeding the highest speed in the unlimited span Open Class, an exceptional result. Masak went on to win the 1993 U.S. 15 Meter Nationals gliding competition, using winglets on his prototype Masak Scimitar. The Masak winglets were originally retrofitted to production sailplanes, but within 10 years of their introduction, most high-performance gliders were equipped from the factory with winglets or other wingtip devices. It took over a decade for winglets to first appear on a production airliner, the original application that was the focus of the NASA development. Yet, once the advantages of winglets were proven in competition, adoption was swift with gliders. The point difference between the winner and the runner-up in soaring competition is often less than one percent, so even a small improvement in efficiency is a significant competitive advantage. Many non-competition pilots fitted winglets for handling benefits such as increased roll rate and roll authority and reduced tendency for wing tip stall. The benefits are notable, because sailplane winglets must be removable to allow the glider to be stored in a trailer, so they are usually installed only at the pilot's preference. The Glaser-Dirks DG-303, an early glider derivative design, incorporating winglets as factory standard equipment. Non-planar wingtip Aviation Partners developed and flight tested a closed-surface Spiroid winglet on a Falcon 50 in 2010. Non-planar wingtips are normally angled upwards in a polyhedral wing configuration, increasing the local dihedral near the wing tip, with polyhedral wing designs themselves having been popular on free-flight model aircraft designs for decades. Non-planar wingtips provide the wake control benefit of winglets, with less parasitic drag penalty, if designed carefully. The non-planar wing tip is often swept back like a raked wingtip and may also be combined with a winglet. A winglet is also a special case of a non-planar wingtip. Aircraft designers employed mostly planar wing designs with simple dihedral after World War II, prior to the introduction of winglets. With the wide acceptance of winglets in new sailplane designs of the 1990s, designers sought to further optimize the aerodynamic performance of their wingtip designs. Glider winglets were originally retrofitted directly to planar wings, with only a small, nearly right-angle, transition area. Once the performance of the winglet itself was optimized, attention was turned to the transition between the wing and winglet. A common application was tapering the transition area from the wing tip chord to the winglet chord and raking the transition area back, to place the winglet in the optimal position. If the tapered portion was canted upward, the winglet height could also be reduced. Eventually, designers employed multiple non-planar sections, each canting up at a greater angle, dispensing with the winglets entirely. The Schempp-Hirth Discus-2 and Schempp-Hirth Duo Discus use non-planar wingtips. Active wingtip device Tamarack Aerospace Group, a company founded in 2010 by aerospace structural engineer Nicholas Guida, has patented an Active Technology Load Alleviation System (ATLAS), a modified version of a wingtip device. The system uses Tamarack Active Camber Surfaces (TACS) to aerodynamically "switch off" the effects of the wingtip device when the aircraft is experiencing high-g events such as large gusts or severe pull-ups. TACS are movable panels, similar to flaps or ailerons, on the trailing edge of the wing extension. The system is controlled by the aircraft's electrical system and a high-speed servo which is activated when the aircraft senses an oncoming stress event, essentially simulating an actuating wingtip. However, the wingtip itself is fixed and the TACS are the only moving part of the wingtip system. Tamarack first introduced ATLAS for the Cessna Citation family aircraft, and it has been certified for use by the Federal Aviation Administration and European Union Aviation Safety Agency. In December 2024, Tamarack Aerospace had installed 200 Active Winglet on CitationJet airplanes. Actuating wingtip device There has been research into actuating wingtip devices, including a filed patent application, though no aircraft currently uses this feature as described. The XB-70 Valkyrie's wingtips were capable of drooping downward in flight, to facilitate Mach 3 flight using waveriding. Use on rotating blades Wingtip devices are also used on rotating propeller, helicopter rotor, and wind turbine blades to reduce drag, reduce diameter, reduce noise and/or improve efficiency. By reducing aircraft blade tip vortices interacting with the ground surface during taxiing, takeoff, and hover, these devices can reduce damage from dirt and small stones picked up in the vortices. Rotorcraft applications The main rotor blades of the AgustaWestland AW101 (formerly the EH101) have a distinctive tip shape; pilots have found that this rotor design alters the downwash field and reduces brownout which limits visibility in dusty areas and leads to accidents. Propeller applications Hartzell Propeller developed their "Q-tip" propeller used on the Piper PA-42 Cheyenne and several other fixed-wing aircraft types by bending the blade tips back at a 90-degree angle to get the same thrust from a reduced diameter propeller disk; the reduced propeller tip speed reduces noise, according to the manufacturer. Modern scimitar propellers have increased sweepback at the tips, resembling a raked tip on an aircraft wing. Other applications Some ceiling fans have wingtip devices. Fan manufacturer Big Ass Fans has claimed that their Isis fan, equipped with wingtip devices, has superior efficiency. However, for certain high-volume, low-speed designs, wingtip devices may not improve efficiency. Another application of the same principle was introduced to the keel of the "America's Cup"- winning Australian yacht Australia II of 1982, designed by Ben Lexcen.	Technology	Aircraft components	null
313418	https://en.wikipedia.org/wiki/Luminous%20intensity	Luminous intensity	In photometry, luminous intensity is a measure of the wavelength-weighted power emitted by a light source in a particular direction per unit solid angle, based on the luminosity function, a standardized model of the sensitivity of the human eye. The SI unit of luminous intensity is the candela (cd), an SI base unit. Measurement Photometry deals with the measurement of visible light as perceived by human eyes. The human eye can only see light in the visible spectrum and has different sensitivities to light of different wavelengths within the spectrum. When adapted for bright conditions (photopic vision), the eye is most sensitive to yellow-green light at 555 nm. Light with the same radiant intensity at other wavelengths has a lower luminous intensity. The curve which represents the response of the human eye to light is a defined standard function or established by the International Commission on Illumination (CIE, for Commission Internationale de l'Éclairage) and standardized in collaboration with the ISO. Luminous intensity of artificial light sources is typically measured using and a goniophotometer outfitted with a photometer or a spectroradiometer. Relationship to other measures Luminous intensity should not be confused with another photometric unit, luminous flux, which is the total perceived power emitted in all directions. Luminous intensity is the perceived power per unit solid angle. If a lamp has a 1 lumen bulb and the optics of the lamp are set up to focus the light evenly into a 1 steradian beam, then the beam would have a luminous intensity of 1 candela. If the optics were changed to concentrate the beam into 1/2 steradian then the source would have a luminous intensity of 2 candela. The resulting beam is narrower and brighter, though its luminous flux remains unchanged. Luminous intensity is also not the same as the radiant intensity, the corresponding objective physical quantity used in the measurement science of radiometry. Units Like other SI base units, the candela has an operational definition—it is defined by the description of a physical process that will produce one candela of luminous intensity. By definition, if one constructs a light source that emits monochromatic green light with a frequency of 540 THz, and that has a radiant intensity of 1/683 watts per steradian in a given direction, that light source will emit one candela in the specified direction. The frequency of light used in the definition corresponds to a wavelength in a vacuum of , which is near the peak of the eye's response to light. If the source emitted uniformly in all directions, the total radiant flux would be about , since there are 4 steradians in a sphere. A typical modern candle produces very roughly one candela while releasing heat at roughly . Prior to the definition of the candela, a variety of units for luminous intensity were used in various countries. These were typically based on the brightness of the flame from a "standard candle" of defined composition, or the brightness of an incandescent filament of specific design. One of the best-known of these standards was the English standard: candlepower. One candlepower was the light produced by a pure spermaceti candle weighing one sixth of a pound and burning at a rate of 120 grains per hour. Germany, Austria, and Scandinavia used the Hefnerkerze, a unit based on the output of a Hefner lamp. In 1881, Jules Violle proposed the Violle as a unit of luminous intensity, and it was notable as the first unit of light intensity that did not depend on the properties of a particular lamp. All of these units were superseded by the definition of the candela. Usage The luminous intensity for monochromatic light of a particular wavelength is given by where is the luminous intensity in candelas (cd), is the radiant intensity in watts per steradian (W/sr), is the standard luminosity function. If more than one wavelength is present (as is usually the case), one must sum or integrate over the spectrum of wavelengths present to get the luminous intensity:	Physical sciences	Optics	Physics
313530	https://en.wikipedia.org/wiki/Sperm%20whale	Sperm whale	The sperm whale or cachalot (Physeter macrocephalus) is the largest of the toothed whales and the largest toothed predator. It is the only living member of the genus Physeter and one of three extant species in the sperm whale family, along with the pygmy sperm whale and dwarf sperm whale of the genus Kogia. The sperm whale is a pelagic mammal with a worldwide range, and will migrate seasonally for feeding and breeding. Females and young males live together in groups, while mature males (bulls) live solitary lives outside of the mating season. The females cooperate to protect and nurse their young. Females give birth every four to twenty years, and care for the calves for more than a decade. A mature, healthy sperm whale has no natural predators, although calves and weakened adults are sometimes killed by pods of killer whales (orcas). Mature males average in length, with the head representing up to one-third of the animal's length. Plunging to , it is the third deepest diving mammal, exceeded only by the southern elephant seal and Cuvier's beaked whale. The sperm whale uses echolocation and vocalization with source level as loud as 236 decibels (re 1 μPa m) underwater, the loudest of any animal. It has the largest brain on Earth, more than five times heavier than a human's. Sperm whales can live 70 years or more. Sperm whales' heads are filled with a waxy substance called "spermaceti" (sperm oil), from which the whale derives its name. Spermaceti was a prime target of the whaling industry and was sought after for use in oil lamps, lubricants, and candles. Ambergris, a solid waxy waste product sometimes present in its digestive system, is still highly valued as a fixative in perfumes, among other uses. Beachcombers look out for ambergris as flotsam. Sperm whaling was a major industry in the 19th century, depicted in the novel Moby-Dick. The species is protected by the International Whaling Commission moratorium, and is listed as vulnerable by the International Union for Conservation of Nature. Taxonomy and naming Etymology The name "sperm whale" is a clipping of "spermaceti whale". Spermaceti, originally mistakenly identified as the whales' semen, is the semi-liquid, waxy substance found within the whale's head. (See "Spermaceti organ and melon" below.) The sperm whale is also known as the "cachalot", which is thought to derive from the archaic French for 'tooth' or 'big teeth', as preserved for example in the word in the Gascon dialect (a word of either Romance or Basque origin). The etymological dictionary of Corominas says the origin is uncertain, but it suggests that it comes from the Vulgar Latin 'sword hilts'. The word cachalot came to English via French from Spanish or Portuguese , perhaps from Galician/Portuguese 'big head'. The term is retained in the Russian word for the animal, (), as well as in many other languages. The scientific genus name Physeter comes from the Greek (), meaning 'blowpipe, blowhole (of a whale)', or – as a pars pro toto – 'whale'. The specific name macrocephalus is Latinized from the Greek ( 'big-headed'), from () + (). Its synonymous specific name catodon means 'down-tooth', from the Greek elements ('below') and ('tooth'); so named because it has visible teeth only in its lower jaw. (See "Jaws and teeth" below.) Another synonym australasianus ('Australasian') was applied to sperm whales in the Southern Hemisphere. Taxonomy The sperm whale belongs to the order Cetartiodactyla, the order containing all cetaceans and even-toed ungulates. It is a member of the unranked clade Cetacea, with all the whales, dolphins, and porpoises, and further classified into Odontoceti, containing all the toothed whales and dolphins. It is the sole extant species of its genus, Physeter, in the family Physeteridae. Two species of the related extant genus Kogia, the pygmy sperm whale Kogia breviceps and the dwarf sperm whale K. sima, are placed either in this family or in the family Kogiidae. In some taxonomic schemes the families Kogiidae and Physeteridae are combined as the superfamily Physeteroidea (see the separate entry on the sperm whale family). Swedish ichthyologist Peter Artedi described it as Physeter catodon in his 1738 work Genera piscium, from the report of a beached specimen in Orkney in 1693 and two beached in the Netherlands in 1598 and 1601. The 1598 specimen was near Berkhey. The sperm whale is one of the species originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae. He recognised four species in the genus Physeter. Experts soon realised that just one such species exists, although there has been debate about whether this should be named P. catodon or P. macrocephalus, two of the names used by Linnaeus. Both names are still used, although most recent authors now accept macrocephalus as the valid name, limiting catodon status to a lesser synonym. Until 1974, the species was generally known as P. catodon. In that year, however, Dutch zoologists Antonius M. Husson and Lipke Holthuis proposed that the correct name should be P. macrocephalus, the second name in the genus Physeter published by Linnaeus concurrently with P. catodon. This proposition was based on the grounds that the names were synonyms published simultaneously, and, therefore, the ICZN Principle of the First Reviser should apply. In this instance, it led to the choice of P. macrocephalus over P. catodon, a view re-stated in Holthuis, 1987. This has been adopted by most subsequent authors, although Schevill (1986 and 1987) argued that macrocephalus was published with an inaccurate description and that therefore only the species catodon was valid, rendering the principle of "First Reviser" inapplicable. The most recent version of ITIS has altered its usage from P. catodon to P. macrocephalus, following L. B. Holthuis and more recent (2008) discussions with relevant experts. Furthermore, The Taxonomy Committee of the Society for Marine Mammalogy, the largest international association of marine mammal scientists in the world, officially uses Physeter macrocephalus when publishing their definitive list of marine mammal species. Biology External appearance The sperm whale is the largest toothed whale and is among the most sexually dimorphic of all cetaceans. Both sexes are about the same size at birth, but mature males are typically 30% to 50% longer and three times as massive as females. Newborn sperm whales are usually between long. Female sperm whales are sexually mature at in length, whilst males are sexually mature at . Female sperm whales are physically mature at about in length and generally do not achieve lengths greater than . The largest female sperm whale measured up to long, and an individual of such size would have weighed about . Male sperm whales are physically mature at about in length, and larger males can generally achieve . An long male sperm whale is estimated to have weighed . By contrast, the second largest toothed whale (Baird's beaked whale) measures up to and weighs up to . There are occasional reports of individual sperm whales achieving even greater lengths, with some historical claims reaching or exceeding . One example is the whale that sank the Essex (one of the incidents behind Moby-Dick), which was claimed to be . However, there is disagreement as to the accuracy of some of these claims, which are often considered exaggerations or as being measured along the curves of the body. An individual measuring was reported from a Soviet whaling fleet near the Kuril Islands in 1950 and is cited by some authors as the largest accurately measured. It has been estimated to weigh . In a review of size variation in marine megafauna, McClain and colleagues noted that the International Whaling Commission's data contained eight individuals larger than . The authors supported a male from the South Pacific in 1933 as the largest recorded. However, sizes like these are rare, with 95% of recorded sperm whales below 15.85 metres (52.0 ft). In 1853, one sperm whale was reported at in length, with a head measuring . Large lower jawbones are held in the British Natural History Museum and the Oxford University Museum of Natural History, measuring and , respectively. The average size of sperm whales has decreased over the years, probably due to pressure from whaling. Another view holds that exploitation by overwhaling had virtually no effect on the size of the bull sperm whales, and their size may have actually increased in current times on the basis of density dependent effects. Old males taken at Solander Islands were recorded to be extremely large and unusually rich in blubbers. The sperm whale's unique body is unlikely to be confused with any other species. The sperm whale's distinctive shape comes from its very large, block-shaped head, which can be one-quarter to one-third of the animal's length. The S-shaped blowhole is located very close to the front of the head and shifted to the whale's left. This gives rise to a distinctive bushy, forward-angled spray. The sperm whale's flukes (tail lobes) are triangular and very thick. Proportionally, they are larger than that of any other cetacean, and are very flexible. The whale lifts its flukes high out of the water as it begins a feeding dive. It has a series of ridges on the back's caudal third instead of a dorsal fin. The largest ridge was called the 'hump' by whalers, and can be mistaken for a dorsal fin because of its shape and size. In contrast to the smooth skin of most large whales, its back skin is usually wrinkly and has been likened to a prune by whale-watching enthusiasts. Albinos have been reported. Skeleton The ribs are bound to the spine by flexible cartilage, which allows the ribcage to collapse rather than snap under high pressure. While sperm whales are well adapted to diving, repeated dives to great depths have long-term effects. Bones show the same avascular necrosis that signals decompression sickness in humans. Older skeletons showed the most extensive damage, whereas calves showed no damage. This damage may indicate that sperm whales are susceptible to decompression sickness, and sudden surfacing could be lethal to them. Like that of all cetaceans, the spine of the sperm whale has reduced zygapophysial joints, of which the remnants are modified and are positioned higher on the vertebral dorsal spinous process, hugging it laterally, to prevent extensive lateral bending and facilitate more dorso-ventral bending. These evolutionary modifications make the spine more flexible but weaker than the spines of terrestrial vertebrates. Like many cetaceans, the sperm whale has a vestigial pelvis that is not connected to the spine. Like that of other toothed whales, the skull of the sperm whale is asymmetrical so as to aid echolocation. Sound waves that strike the whale from different directions will not be channeled in the same way. Within the basin of the cranium, the openings of the bony narial tubes (from which the nasal passages spring) are skewed towards the left side of the skull. Jaws and teeth The sperm whale's lower jaw is very narrow and underslung. The sperm whale has 18 to 26 teeth on each side of its lower jaw which fit into sockets in the upper jaw. The teeth are cone-shaped and weigh up to each. The teeth are functional, but do not appear to be necessary for capturing or eating squid, as well-fed animals have been found without teeth or even with deformed jaws. One hypothesis is that the teeth are used in aggression between males. Mature males often show scars which seem to be caused by the teeth. Rudimentary teeth are also present in the upper jaw, but these rarely emerge into the mouth. Analyzing the teeth is the preferred method for determining a whale's age. Like the age-rings in a tree, the teeth build distinct layers of cementum and dentine as they grow. Brain The sperm whale brain is the largest known of any modern or extinct animal, weighing on average about (with the smallest known weighing and the largest known weighing ), more than five times heavier than a human brain, and has a volume of about 8,000 cm3. Although larger brains generally correlate with higher intelligence, it is not the only factor. Elephants and dolphins also have larger brains than humans. The sperm whale has a lower encephalization quotient than many other whale and dolphin species, lower than that of non-human anthropoid apes, and much lower than that of humans. The sperm whale's cerebrum is the largest in all mammalia, both in absolute and relative terms. The olfactory system is reduced, suggesting that the sperm whale has a poor sense of taste and smell. By contrast, the auditory system is enlarged. The pyramidal tract is poorly developed, reflecting the reduction of its limbs. Biological systems The sperm whale respiratory system has adapted to cope with drastic pressure changes when diving. The flexible ribcage allows lung collapse, reducing nitrogen intake, and metabolism can decrease to conserve oxygen. Between dives, the sperm whale surfaces to breathe for about eight minutes before diving again. Odontoceti (toothed whales) breathe air at the surface through a single, S-shaped blowhole, which is extremely skewed to the left. Sperm whales spout (breathe) 3–5 times per minute at rest, increasing to 6–7 times per minute after a dive. The blow is a noisy, single stream that rises up to or more above the surface and points forward and left at a 45° angle. On average, females and juveniles blow every 12.5 seconds before dives, while large males blow every 17.5 seconds before dives. A sperm whale killed south of Durban, South Africa, after a 1-hour, 50-minute dive was found with two dogfish (Scymnodon sp.), usually found at the sea floor, in its belly. The sperm whale has the longest intestinal system in the world, exceeding 300 m in larger specimens. The sperm whale has a four-chambered stomach that is similar to ruminants. The first secretes no gastric juices and has very thick muscular walls to crush the food (since whales cannot chew) and resist the claw and sucker attacks of swallowed squid. The second chamber is larger and is where digestion takes place. Undigested squid beaks accumulate in the second chamber – as many as 18,000 have been found in some dissected specimens. Most squid beaks are vomited by the whale, but some occasionally make it to the hindgut. Such beaks precipitate the formation of ambergris. In 1959, the heart of a 22 metric-ton (24 short-ton) male taken by whalers was measured to be , about 0.5% of its total mass. The circulatory system has a number of specific adaptations for the aquatic environment. The diameter of the aortic arch increases as it leaves the heart. This bulbous expansion acts as a windkessel, ensuring a steady blood flow as the heart rate slows during diving. The arteries that leave the aortic arch are positioned symmetrically. There is no costocervical artery. There is no direct connection between the internal carotid artery and the vessels of the brain. Their circulatory system has adapted to dive at great depths, as much as for up to 120 minutes. More typical dives are around and 35 minutes in duration. Myoglobin, which stores oxygen in muscle tissue, is much more abundant than in terrestrial animals. The blood has a high density of red blood cells, which contain oxygen-carrying haemoglobin. The oxygenated blood can be directed towards only the brain and other essential organs when oxygen levels deplete. The spermaceti organ may also play a role by adjusting buoyancy (see below). The arterial retia mirabilia are extraordinarily well-developed. The complex arterial retia mirabilia of the sperm whale are more extensive and larger than those of any other cetacean. Senses Spermaceti organ and melon Atop the whale's skull is positioned a large complex of organs filled with a liquid mixture of fats and waxes called spermaceti. The purpose of this complex is to generate powerful and focused clicking sounds, the existence of which was proven by Valentine Worthington and William Schevill when a recording was produced on a research vessel in May 1959. The sperm whale uses these sounds for echolocation and communication. The spermaceti organ is like a large barrel of spermaceti. Its surrounding wall, known as the case, is extremely tough and fibrous. The case can hold within it up to 1,900 litres of spermaceti. It is proportionately larger in males. This oil is a mixture of triglycerides and wax esters. It has been suggested that it is homologous to the dorsal bursa organ found in dolphins. The proportion of wax esters in the spermaceti organ increases with the age of the whale: 38–51% in calves, 58–87% in adult females, and 71–94% in adult males. The spermaceti at the core of the organ has a higher wax content than the outer areas. The speed of sound in spermaceti is 2,684 m/s (at 40 kHz, 36 °C), making it nearly twice as fast as in the oil in a dolphin's melon. Below the spermaceti organ lies the "junk" which consists of compartments of spermaceti separated by cartilage. It is analogous to the melon found in other toothed whales. The structure of the junk redistributes physical stress across the skull and may have evolved to protect the head during ramming. Running through the head are two air passages. The left passage runs alongside the spermaceti organ and goes directly to the blowhole, whilst the right passage runs underneath the spermaceti organ and passes air through a pair of phonic lips and into the distal sac at the very front of the nose. The distal sac is connected to the blowhole and the terminus of the left passage. When the whale is submerged, it can close the blowhole, and air that passes through the phonic lips can circulate back to the lungs. The sperm whale, unlike other odontocetes, has only one pair of phonic lips, whereas all other toothed whales have two, and it is located at the front of the nose instead of behind the melon. At the posterior end of this spermaceti complex is the frontal sac, which covers the concave surface of the cranium. The posterior wall of the frontal sac is covered with fluid-filled knobs, which are about 4–13 mm in diameter and separated by narrow grooves. The anterior wall is smooth. The knobbly surface reflects sound waves that come through the spermaceti organ from the phonic lips. The grooves between the knobs trap a film of air that is consistent whatever the orientation or depth of the whale, making it an excellent sound mirror. The spermaceti organs may also help adjust the whale's buoyancy. It is hypothesized that before the whale dives, cold water enters the organ, and it is likely that the blood vessels constrict, reducing blood flow, and, hence, temperature. The wax therefore solidifies and reduces in volume. The increase in specific density generates a down force of about and allows the whale to dive with less effort. During the hunt, oxygen consumption, together with blood vessel dilation, produces heat and melts the spermaceti, increasing its buoyancy and enabling easy surfacing. However, more recent work has found many problems with this theory including the lack of anatomical structures for the actual heat exchange. Another issue is that if the spermaceti does indeed cool and solidify, it would affect the whale's echolocation ability just when it needs it to hunt in the depths. Herman Melville's fictional story Moby-Dick suggests that the "case" containing the spermaceti serves as a battering ram for use in fights between males. A few famous instances include the well-documented sinking of the ships Essex and Ann Alexander by attackers estimated to weigh only one-fifth as much as the ships. Eyes and vision The sperm whale's eye does not differ greatly from those of other toothed whales except in size. It is the largest among the toothed whales, weighing about 170 g. It is overall ellipsoid in shape, compressed along the visual axis, measuring about 7×7×3 cm. The cornea is elliptical and the lens is spherical. The sclera is very hard and thick, roughly 1 cm anteriorly and 3 cm posteriorly. There are no ciliary muscles. The choroid is very thick and contains a fibrous tapetum lucidum. Like other toothed whales, the sperm whale can retract and protrude its eyes, thanks to a 2-cm-thick retractor muscle attached around the eye at the equator, but are unable to roll the eyes in their sockets. According to Fristrup and Harbison (2002), sperm whale's eyes afford good vision and sensitivity to light. They conjectured that sperm whales use vision to hunt squid, either by detecting silhouettes from below or by detecting bioluminescence. If sperm whales detect silhouettes, Fristrup and Harbison suggested that they hunt upside down, allowing them to use the forward parts of the ventral visual fields for binocular vision. Sleeping For some time researchers have been aware that pods of sperm whales may sleep for short periods, assuming a vertical position with their heads just below or at the surface, or head down. A 2008 study published in Current Biology recorded evidence that whales may sleep with both sides of the brain. It appears that some whales may fall into a deep sleep for about 7 percent of the time, most often between 6 p.m. and midnight. Genetics Sperm whales have 21 pairs of chromosomes (2n=42). The genome of live whales can be examined by recovering shed skin. Vocalization complex After Valentine Worthington and William E. Schevill confirmed the existence of sperm whale vocalization, further studies found that sperm whales are capable of emitting sounds at a source level of 230 decibels–making the sperm whale the loudest animal in the world. Mechanism When echolocating, the sperm whale emits a directionally focused beam of broadband clicks. Clicks are generated by forcing air through a pair of phonic lips (also known as "monkey lips" or "") at the front end of the nose, just below the blowhole. The sound then travels backwards along the length of the nose through the spermaceti organ. Most of the sound energy is then reflected off the frontal sac at the cranium and into the melon, whose lens-like structure focuses it. Some of the sound will reflect back into the spermaceti organ and back towards the front of the whale's nose, where it will be reflected through the spermaceti organ a third time. This back and forth reflection which happens on the scale of a few milliseconds creates a multi-pulse click structure. This multi-pulse click structure allows researchers to measure the whale's spermaceti organ using only the sound of its clicks. Because the interval between pulses of a sperm whale's click is related to the length of the sound producing organ, an individual whale's click is unique to that individual. However, if the whale matures and the size of the spermaceti organ increases, the tone of the whale's click will also change. The lower jaw is the primary reception path for the echoes. A continuous fat-filled canal transmits received sounds to the inner ear. The source of the air forced through the phonic lips is the right nasal passage. While the left nasal passage opens to the blow hole, the right nasal passage has evolved to supply air to the phonic lips. It is thought that the nostrils of the land-based ancestor of the sperm whale migrated through evolution to their current functions, the left nostril becoming the blowhole and the right nostril becoming the phonic lips. Air that passes through the phonic lips passes into the distal sac, then back down through the left nasal passage. This recycling of air allows the whale to continuously generate clicks for as long as it is submerged. Vocalization types The sperm whale's vocalizations are all based on clicking, described in four types: the usual echolocation, creaks, codas, and slow clicks. The usual echolocation click type is used in searching for prey. A creak is a rapid series of high-frequency clicks that sounds somewhat like a creaky door hinge. It is typically used when homing in on prey. Slow clicks are heard only in the presence of males (it is not certain whether females occasionally make them). Males make a lot of slow clicks in breeding grounds (74% of the time), both near the surface and at depth, which suggests they are primarily mating signals. Outside breeding grounds, slow clicks are rarely heard, and usually near the surface. Codas The most distinctive vocalizations are codas, which are short rhythmic sequences of clicks, mostly numbering 3–12 clicks, in stereotyped patterns. They are classified using variations in the number of clicks, rhythm, and tempo. Codas are the result of vocal learning within a stable social group, and are made in the context of the whales' social unit. "The foundation of sperm whale society is the matrilineally based social unit of ten or so females and their offspring. The members of the unit travel together, suckle each others' infants, and babysit them while mothers make long deep dives to feed." Over 70% of a sperm whale's time is spent independently foraging; codas "could help whales reunite and reaffirm their social ties in between long foraging dives." While nonidentity codas are commonly used in multiple different clans, some codas express clan identity, and denote different patterns of travel, foraging, and socializing or avoidance among clans. In particular, whales will not group with whales of another clan even though they share the same geographical area. Statistically, as the clans' ranges become more overlapped, the distinction in clan identity coda usage becomes more pronounced. Distinctive codas identify seven clans described among the approximately 150,000 female sperm whales in the Pacific Ocean, and there are another four clans in the Atlantic. As "arbitrary traits that function as reliable indicators of cultural group membership," clan identity codas act as symbolic markers that modulate interactions between individuals. Individual identity in sperm whale vocalizations is an ongoing scientific issue, however. A distinction needs to be made between cues and signals. Human acoustic tools can distinguish individual whales by analyzing micro-characteristics of their vocalizations, and the whales can probably do the same. This does not prove that the whales deliberately use some vocalizations to signal individual identity in the manner of the signature whistles that bottlenose dolphins use as individual labels. Ecology Distribution Sperm whales are among the most cosmopolitan species. They prefer ice-free waters over deep. Although both sexes range through temperate and tropical oceans and seas, only adult males populate higher latitudes. Among several regions, such as along coastal waters of southern Australia, sperm whales have been considered to be locally extinct. They are relatively abundant from the poles to the equator and are found in all the oceans. They inhabit the Mediterranean Sea, but not the Black Sea, while their presence in the Red Sea is uncertain. The shallow entrances to both the Black Sea and the Red Sea may account for their absence. The Black Sea's lower layers are also anoxic and contain high concentrations of sulphur compounds such as hydrogen sulphide. The first ever sighting off the coast of Pakistan was made in 2017. The first ever record off the west coast of the Korean Peninsula (Yellow Sea) was made in 2005. followed by one near Ganghwa Island in 2009. Populations are denser close to continental shelves and canyons. Sperm whales are usually found in deep, off-shore waters, but may be seen closer to shore, in areas where the continental shelf is small and drops quickly to depths of . Coastal areas with significant sperm whale populations include the Azores and Dominica. In east Asian waters, whales are also observed regularly in coastal waters in places such as the Commander and Kuril Islands, Shiretoko Peninsula which is one of few locations where sperm whales can be observed from shores, off Kinkasan, vicinity to Tokyo Bay and the Bōsō Peninsula to the Izu and the Izu Islands, the Volcano Islands, Yakushima and the Tokara Islands to the Ryukyu Islands, Taiwan, the Northern Mariana Islands, and so forth. Historical catch records suggest there could have been smaller aggression grounds in the Sea of Japan as well. Along the Korean Peninsula, the first confirmed observation within the Sea of Japan, eight animals off Guryongpo, was made in 2004 since after the last catches of five whales off Ulsan in 1911, while nine whales were observed in the East China Sea side of the peninsula in 1999. Grown males are known to enter surprisingly shallow bays to rest (whales will be in a state of rest during these occasions). Unique, coastal groups have been reported from various areas around the globe, such as near Scotland's coastal waters, and the Shiretoko Peninsula, off Kaikōura, in Davao Gulf. Such coastal groups were more abundant in pre-whaling days. Genetic analysis indicates that the world population of sperm whales originated in the Pacific Ocean from a population of about 10,000 animals around 100,000 years ago, when expanding ice caps blocked off their access to other seas. In particular, colonization of the Atlantic was revealed to have occurred multiple times during this expansion of their range. Diet Sperm whales usually dive between , and sometimes , in search of food. Such dives can last more than an hour. They feed on several species, notably the giant squid, but also the colossal squid, octopuses, and fish such as demersal rays and sharks, but their diet is mainly medium-sized squid. Sperm whales may also possibly prey upon swordfish on rare occasions. Some prey may be taken accidentally while eating other items. Most of what is known about deep-sea squid has been learned from specimens in captured sperm whale stomachs, although more recent studies analysed faeces. One study, carried out around the Galápagos, found that squid from the genera Histioteuthis (62%), Ancistrocheirus (16%), and Octopoteuthis (7%) weighing between were the most commonly taken. Battles between sperm whales and giant squid or colossal squid have never been observed by humans; however, white scars are believed to be caused by the large squid. One study published in 2010 collected evidence that suggests that female sperm whales may collaborate when hunting Humboldt squid. Tagging studies have shown that sperm whales hunt upside down at the bottom of their deep dives. It is suggested that the whales can see the squid silhouetted above them against the dim surface light. An older study, examining whales captured by the New Zealand whaling fleet in the Cook Strait region, found a 1.69:1 ratio of squid to fish by weight. Sperm whales sometimes take sablefish and toothfish from long lines. Long-line fishing operations in the Gulf of Alaska complain that sperm whales take advantage of their fishing operations to eat desirable species straight off the line, sparing the whales the need to hunt. However, the amount of fish taken is very little compared to what the sperm whale needs per day. Video footage has been captured of a large male sperm whale "bouncing" a long line, to gain the fish. Sperm whales are believed to prey on the megamouth shark, a rare and large deep-sea species discovered in the 1970s. In one case, three sperm whales were observed attacking or playing with a megamouth. Sperm whales have also been noted to feed on bioluminescent pyrosomes such as Pyrosoma atlanticum. It is thought that the foraging strategy of sperm whales for bioluminescent squids may also explain the presence of these light-emitting pyrosomes in the diet of the sperm whale. The sharp beak of a consumed squid lodged in the whale's intestine may lead to the production of ambergris, analogous to the production of pearls in oysters. The irritation of the intestines caused by squid beaks stimulates the secretion of this lubricant-like substance. Sperm whales are prodigious feeders and eat around 3% of their body weight per day. The total annual consumption of prey by sperm whales worldwide is estimated to be about . In comparison, human consumption of seafood is estimated to be . Sperm whales hunt through echolocation. Their clicks are among the most powerful sounds in the animal kingdom (see above). It has been hypothesised that it can stun prey with its clicks. Experimental studies attempting to duplicate this effect have been unable to replicate the supposed injuries, casting doubt on this idea. One study showing that sound pressure levels on the squid are more than an order of magnitude below levels required for debilitation, and therefore, precluding acoustic stunning to facilitate prey capture. Sperm whales, as well as other large cetaceans, help fertilise the surface of the ocean by consuming nutrients in the depths and transporting those nutrients to the oceans' surface when they defecate, an effect known as the whale pump. This fertilises phytoplankton and other plants on the surface of the ocean and contributes to ocean productivity and the drawdown of atmospheric carbon. Life cycle Sperm whales can live 70 years or more. They are a prime example of a species that has been K-selected, meaning their reproductive strategy is associated with stable environmental conditions and comprises a low birth rate, significant parental aid to offspring, slow maturation, and high longevity. How they choose mates has not been definitively determined. Bulls will fight with each other over females, and males will mate with multiple females, making them polygynous, but they do not dominate the group as in a harem. Bulls do not provide paternal care to their offspring but rather play a fatherly role to younger bulls to show dominance. Females become fertile at around 9 years of age. The oldest pregnant female ever recorded was 41 years old. Gestation requires 14 to 16 months, producing a single calf. Sexually mature females give birth once every 4 to 20 years (pregnancy rates were higher during the whaling era). Birth is a social event, as the mother and calf need others to protect them from predators. The other adults may jostle and bite the newborn in its first hours. Lactation proceeds for 19 to 42 months, but calves, rarely, may suckle up to 13 years. Like that of other whales, the sperm whale's milk has a higher fat content than that of terrestrial mammals: about 36%, compared to 4% in cow milk. This gives it a consistency similar to cottage cheese, which prevents it from dissolving in the water before the calf can drink it. It has an energy content of roughly 3,840 kcal/kg, compared to just 640 kcal/kg in cow milk. Calves may be allowed to suckle from females other than their mothers. Males become sexually mature at 18 years. Upon reaching sexual maturity, males move to higher latitudes, where the water is colder and feeding is more productive. Females remain at lower latitudes. Males reach their full size at about age 50. Social behaviour Relations within the species Like elephants, females and their young live in matriarchal groups called pods, while bulls live apart. Bulls sometimes form loose bachelor groups with other males of similar age and size. As they grow older, they typically live solitary lives, only returning to the pod to socialize or to breed. Bulls have beached themselves together, suggesting a degree of cooperation which is not yet fully understood. The whales rarely, if ever, leave their group. A social unit is a group of sperm whales who live and travel together over a period of years. Individuals rarely, if ever, join or leave a social unit. There is a huge variance in the size of social units. They are most commonly between six and nine individuals in size but can have more than twenty. Unlike orcas, sperm whales within a social unit show no significant tendency to associate with their genetic relatives. Females and calves spend about three-quarters of their time foraging and a quarter of their time socializing. Socializing usually takes place in the afternoon. When sperm whales socialize, they emit complex patterns of clicks called codas. They will spend much of the time rubbing against each other. Tracking of diving whales suggests that groups engage in herding of prey, similar to bait balls created by other species, though the research needs to be confirmed by tracking the prey. Relations with other species The most common natural predator of sperm whales is the orca (killer whale), but pilot whales and false killer whales sometimes harass them. Orcas prey on target groups of females with young, usually making an effort to extract and kill a calf. The females will protect their calves or an injured adult by encircling them. They may face inwards with their tails out (the 'marguerite formation', named after the flower). The heavy and powerful tail of an adult whale is potentially capable of delivering lethal blows. Alternatively, they may face outwards (the 'heads-out formation'). Other than sperm whales, southern right whales had been observed to perform similar formations. However, formations in non-dangerous situations have been recorded as well. Early whalers exploited this behaviour, attracting a whole unit by injuring one of its members. Such a tactic is described in Moby-Dick: "Say you strike a Forty-barrel-bull—poor devil! all his comrades quit him. But strike a member of the harem school, and her companions swim around her with every token of concern, sometimes lingering so near her and so long, as themselves to fall a prey."If the killer whale pod is large, its members may sometimes be able to kill adult female sperm whales and can at least injure an entire pod of sperm whales. Bulls have no predators, and are believed to be too large, powerful and aggressive to be threatened by killer whales. Solitary bulls are known to interfere and come to the aid of vulnerable groups nearby. However, the bull sperm whale, when accompanying pods of female sperm whales and their calves as such, may be reportedly unable to effectively dissuade killer whales from their attacks on the group, although the killer whales may end the attack sooner when a bull is present. However, male sperm whales have been observed to attack and intimidate killer whale pods in competitive feeding instances. An incident was filmed from a long-line trawler: a killer whale pod was systematically taking fish caught on the trawler's long lines (as the lines were being pulled into the ship) when a male sperm whale appeared to repeatedly charge the killer whale pod in an attempt to drive them away; it was speculated by the film crew that the sperm whale was attempting to access the same fish. The killer whales employed a tail outward and tail-slapping defensive position against the bull sperm whale similar to that used by female sperm whales against attacking killer whales. However, at some potential feeding sites, the killer whales may prevail over sperm whales even when outnumbered by the sperm whales. Some authors consider the killer whales "usually" behaviorally dominant over sperm whales but express that the two species are "fairly evenly matched", with the killer whales' greater aggression, more considerable biting force for their size and predatory prowess more than compensating for their smaller size. Sperm whales are not known for forging bonds with other species, but it was observed that a bottlenose dolphin with a spinal deformity had been accepted into a pod of sperm whales. They are known to swim alongside other cetaceans such as humpback, fin, minke, pilot, and killer whales on occasion. Parasites Sperm whales can suffer from parasites. Out of 35 sperm whales caught during the 1976–1977 Antarctic whaling season, all of them were infected by Anisakis physeteris (in their stomachs) and Phyllobothrium delphini (in their blubber). Both whales with a placenta were infected with Placentonema gigantissima, potentially the largest nematode worm ever described. Evolutionary history Fossil record Although the fossil record is poor, several extinct genera have been assigned to the clade Physeteroidea, which includes the last common ancestor of the modern sperm whale, pygmy sperm whales, dwarf sperm whales, and extinct physeteroids. These fossils include Ferecetotherium, Idiorophus, Diaphorocetus, Aulophyseter, Orycterocetus, Scaldicetus, Placoziphius, Zygophyseter and Acrophyseter. Ferecetotherium, found in Azerbaijan and dated to the late Oligocene (about ), is the most primitive fossil that has been found, which possesses sperm whale-specific features, such as an asymmetric rostrum ("beak" or "snout"). Most sperm whale fossils date from the Miocene period, . Diaphorocetus, from Argentina, has been dated to the early Miocene. Fossil sperm whales from the Middle Miocene include Aulophyseter, Idiorophus and Orycterocetus, all of which were found on the West Coast of the United States, and Scaldicetus, found in Europe and Japan. Orycterocetus fossils have also been found in the North Atlantic Ocean and the Mediterranean Sea, in addition to the west coast of the United States. Placoziphius, found in Europe, and Acrophyseter, from Peru, are dated to the late Miocene. Fossil sperm whales differ from modern sperm whales in tooth count and the shape of the face and jaws. For example, Scaldicetus had a tapered rostrum. Genera from the Oligocene and early and middle Miocene, with the possible exception of Aulophyseter, had teeth in their upper jaws. Acrophyseter, from the late Miocene, also had teeth in both the upper and lower jaws as well as a short rostrum and an upward curving mandible (lower jaw). These anatomical differences suggest that fossil species may not have necessarily been deep-sea squid eaters such as the modern sperm whale, but that some genera mainly ate fish. Zygophyseter, dated from the middle to late Miocene and found in southern Italy, had teeth in both jaws and appears to have been adapted to feed on large prey, rather like the modern killer whale (orca). Other fossil sperm whales with adaptations similar to this are collectively known as killer sperm whales. Two poorly known fossil species belonging to the modern genus Physeter have been recognized so far: P. antiquus (Neogene of France) and P. vetus (Neogene of eastern North America). Physeter vetus is very likely an invalid species, as the few teeth that were used to identify this species appear to be identical to those of another toothed whale, Orycterocetus quadratidens. Phylogeny The traditional view has been that Mysticeti (baleen whales) and Odontoceti (toothed whales) arose from more primitive whales early in the Oligocene period, and that the super-family Physeteroidea, which contains the sperm whale, dwarf sperm whale, and pygmy sperm whale, diverged from other toothed whales soon after that, over . From 1993 to 1996, molecular phylogenetics analyses by Milinkovitch and colleagues, based on comparing the genes of various modern whales, suggested that the sperm whales are more closely related to the baleen whales than they are to other toothed whales, which would have meant that Odontoceti were not monophyletic; in other words, it did not consist of a single ancestral toothed whale species and all its descendants. However, more recent studies, based on various combinations of comparative anatomy and molecular phylogenetics, criticised Milinkovitch's analysis on technical grounds and reaffirmed that the Odontoceti are monophyletic. These analyses also confirm that there was a rapid evolutionary radiation (diversification) of the Physeteroidea in the Miocene period. The Kogiidae (dwarf and pygmy sperm whales) diverged from the Physeteridae (true sperm whales) at least . Usage by humans Sperm whaling Spermaceti, obtained primarily from the spermaceti organ, and sperm oil, obtained primarily from the blubber in the body, were much sought after by 18th, 19th, and 20th century whalers. These substances found a variety of commercial applications, such as candles, soap, cosmetics, machine oil, other specialised lubricants, lamp oil, pencils, crayons, leather waterproofing, rust-proofing materials and many pharmaceutical compounds. Ambergris, a highly expensive, solid, waxy, flammable substance produced in the digestive system of sperm whales, was also sought as a fixative in perfumery. Prior to the early eighteenth century, hunting was mostly by indigenous Indonesians. Legend has it that sometime in the early 18th century, around 1712, Captain Christopher Hussey, while cruising for right whales near shore, was blown offshore by a northerly wind, where he encountered a sperm whale pod and killed one. Although the story may not be true, sperm whales were indeed soon exploited by American whalers. Judge Paul Dudley, in his Essay upon the Natural History of Whales (1725), states that a certain Atkins, 10 or 12 years in the trade, was among the first to catch sperm whales sometime around 1720 off the New England coast. There were only a few recorded instances during the first few decades (1709–1730s) of offshore sperm whaling. Instead, sloops concentrated on the Nantucket Shoals, where they would have taken right whales or went to the Davis Strait region to catch bowhead whales. By the early 1740s, with the advent of spermaceti candles (before 1743), American vessels began to focus on sperm whales. The diary of Benjamin Bangs (1721–1769) shows that, along with the bumpkin sloop he sailed, he found three other sloops flensing sperm whales off the coast of North Carolina in late May 1743. On returning to Nantucket in the summer 1744 on a subsequent voyage, he noted that "45 spermacetes are brought in here this day," another indication that American sperm whaling was in full swing. American sperm whaling soon spread from the east coast of the American colonies to the Gulf Stream, the Grand Banks, West Africa (1763), the Azores (1765), and the South Atlantic (1770s). From 1770 to 1775 Massachusetts, New York, Connecticut, and Rhode Island ports produced 45,000 barrels of sperm oil annually, compared to 8,500 of whale oil. In the same decade, the British began sperm whaling, employing American ships and personnel. By the following decade, the French had entered the trade, also employing American expertise. Sperm whaling increased until the mid-nineteenth century. Spermaceti oil was important in public lighting (for example, in lighthouses, where it was used in the United States until 1862, when it was replaced by lard oil, in turn replaced by petroleum) and for lubricating the machines (such as those used in cotton mills) of the Industrial Revolution. Sperm whaling declined in the second half of the nineteenth century, as petroleum came into broader use. In that sense, petroleum use may be said to have protected whale populations from even greater exploitation. Sperm whaling in the 18th century began with small sloops carrying only one or two whaleboats. The fleet's scope and size increased over time, and larger ships entered the fishery. In the late 18th century and early 19th century, sperm whaling ships sailed to the equatorial Pacific, the Indian Ocean, Japan, the coast of Arabia, Australia and New Zealand. Hunting could be dangerous to the crew, since sperm whales (especially bulls) will readily fight to defend themselves against attack, unlike most baleen whales. When dealing with a threat, sperm whales will use their huge head effectively as a battering ram. Arguably the most famous sperm whale counter-attack occurred on 20 November 1820, when a whale claimed to be about long rammed and sank the Nantucket whaleship Essex. Only 8 out of 21 sailors survived to be rescued by other ships. The sperm whale's ivory-like teeth were often sought by 18th- and 19th-century whalers, who used them to produce inked carvings known as scrimshaw. 30 teeth of the sperm whale can be used for ivory. Each of these teeth, up to and across, are hollow for the first half of their length. Like walrus ivory, sperm whale ivory has two distinct layers. However, sperm whale ivory contains a much thicker inner layer. Though a widely practised art in the 19th century, scrimshaw using genuine sperm whale ivory declined substantially after the retirement of the whaling fleets in the 1880s. Modern whaling was more efficient than open-boat whaling, employing steam-powered ships and exploding harpoons. Initially, modern whaling activity focused on large baleen whales, but as these populations were taken, sperm whaling increased. Spermaceti, the fine waxy oil produced by sperm whales, was in high demand. In both the 1941–1942 and 1942–1943 seasons, Norwegian expeditions took over 3,000 sperm whales off the coast of Peru alone. After World War II, whaling continued unabated to obtain oil for cosmetics and high-performance machinery, such as automobile transmissions. The hunting led to the near-extinction of large whales, including sperm whales, until bans on whale oil use were instituted in 1972. The International Whaling Commission gave the species full protection in 1985, but hunting by Japan in the northern Pacific Ocean continued until 1988. It is estimated that the historic worldwide population numbered 1,100,000 before commercial sperm whaling began in the early 18th century. By 1880, it had declined by an estimated 29 percent. From that date until 1946, the population appears to have partially recovered as whaling activity decreased, but after the Second World War, the population declined even further, to 33 per cent of the pre-whaling population. Between 184,000 and 236,000 sperm whales were killed by the various whaling nations in the 19th century, while in the 20th century, at least 770,000 were taken, the majority between 1946 and 1980. Sperm whales increase levels of primary production and carbon export by depositing iron-rich faeces into surface waters of the Southern Ocean. The iron-rich faeces cause phytoplankton to grow and take up more carbon from the atmosphere. When the phytoplankton dies, it sinks to the deep ocean and takes the atmospheric carbon with it. By reducing the abundance of sperm whales in the Southern Ocean, whaling has resulted in an extra 2 million tonnes of carbon remaining in the atmosphere each year. Remaining sperm whale populations are large enough that the species' conservation status is rated as vulnerable rather than endangered. However, the recovery from centuries of commercial whaling is a slow process, particularly in the South Pacific, where the toll on breeding-age males was severe. Conservation status The total number of sperm whales in the world is unknown, but is thought to be in the hundreds of thousands. The conservation outlook is brighter than for many other whales. Commercial whaling has ceased, and the species is protected almost worldwide, though records indicate that in the 11-year period starting from 2000, Japanese vessels have caught 51 sperm whales. Fishermen do not target sperm whales to eat, but long-line fishing operations in the Gulf of Alaska have complained about sperm whales "stealing" fish from their lines. Since the 2000s , entanglement in fishing nets and collisions with ships represent the greatest threats to the sperm whale population. Other threats include ingestion of marine debris, ocean noise, and chemical pollution. The International Union for Conservation of Nature (IUCN) regards the sperm whale as being "vulnerable". The species is listed as endangered on the United States Endangered Species Act. Sperm whales are listed on Appendix I and Appendix II of the Convention on the Conservation of Migratory Species of Wild Animals (CMS). It is listed on Appendix I as this species has been categorized as being in danger of extinction throughout all or a significant proportion of their range and CMS Parties strive towards strictly protecting these animals, conserving or restoring the places where they live, mitigating obstacles to migration and controlling other factors that might endanger them. It is listed on Appendix II as it has an unfavourable conservation status or would benefit significantly from international co-operation organised by tailored agreements. It is also covered by the Agreement on the Conservation of Cetaceans in the Black Sea, Mediterranean Sea and Contiguous Atlantic Area (ACCOBAMS) and the Memorandum of Understanding for the Conservation of Cetaceans and Their Habitats in the Pacific Islands Region (Pacific Cetaceans MOU). The species is protected under Appendix I of the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES). This makes commercial international trade (including in parts and derivatives) prohibited, with all other international trade strictly regulated through a system of permits and certificates. Cultural importance Rope-mounted teeth are important cultural objects throughout the Pacific. In New Zealand, the Māori know them as "rei puta"; such whale tooth pendants were rare objects because sperm whales were not actively hunted in traditional Māori society. Whale ivory and bone were taken from beached whales. In Fiji the teeth are known as tabua, traditionally given as gifts for atonement or esteem (called sevusevu), and were important in negotiations between rival chiefs. Friedrich Ratzel in The History of Mankind reported in 1896 that, in Fiji, whales' or cachalots' teeth were the most-demanded article of ornament or value. They occurred often in necklaces. Today the tabua remains an important item in Fijian life. The teeth were originally rare in Fiji and Tonga, which exported teeth, but with the Europeans' arrival, teeth flooded the market and this "currency" collapsed. The oversupply led in turn to the development of the European art of scrimshaw. Herman Melville's novel Moby-Dick is based on a true story about a sperm whale that attacked and sank the whaleship Essex. Melville associated the sperm whale with the Bible's Leviathan. The fearsome reputation perpetuated by Melville was based on bull whales' ability to fiercely defend themselves from attacks by early whalers, smashing whaling boats and, occasionally, attacking and destroying whaling ships. In Jules Verne's 1870 novel Twenty Thousand Leagues Under the Seas, the Nautilus fights a group of "cachalots" (sperm whales) to protect a pod of southern right whales from their attacks. Verne portrays them as being savage hunters ("nothing but mouth and teeth"). The sperm whale was designated as the Connecticut state animal by the General Assembly in 1975. It was selected because of its specific contribution to the state's history and because of its present-day plight as an endangered species. Watching sperm whales Sperm whales are not the easiest of whales to watch, due to their long dive times and ability to travel long distances underwater. However, due to the distinctive look and large size of the whale, watching is increasingly popular. Sperm whale watchers often use hydrophones to listen to the clicks of the whales and locate them before they surface. Popular locations for sperm whale watching include the town of Kaikōura on New Zealand's South Island, Andenes and Tromsø in Arctic Norway; as well as the Azores, where the continental shelf is so narrow that whales can be observed from the shore, and Dominica where a long-term scientific research program, The Dominica Sperm Whale Project, has been in operation since 2005. Plastic waste The introduction of plastic waste to the ocean environment by humans is relatively new. From the 1970s, sperm whales have occasionally been found with pieces of plastic in their stomachs.	Biology and health sciences	Cetaceans	null

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