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152547	https://en.wikipedia.org/wiki/Bisection	Bisection	In geometry, bisection is the division of something into two equal or congruent parts (having the same shape and size). Usually it involves a bisecting line, also called a bisector. The most often considered types of bisectors are the segment bisector, a line that passes through the midpoint of a given segment, and the angle bisector, a line that passes through the apex of an angle (that divides it into two equal angles). In three-dimensional space, bisection is usually done by a bisecting plane, also called the bisector. Perpendicular line segment bisector Definition The perpendicular bisector of a line segment is a line which meets the segment at its midpoint perpendicularly. The perpendicular bisector of a line segment also has the property that each of its points is equidistant from segment AB's endpoints: (D). The proof follows from and Pythagoras' theorem: Property (D) is usually used for the construction of a perpendicular bisector: Construction by straight edge and compass In classical geometry, the bisection is a simple compass and straightedge construction, whose possibility depends on the ability to draw arcs of equal radii and different centers: The segment is bisected by drawing intersecting circles of equal radius , whose centers are the endpoints of the segment. The line determined by the points of intersection of the two circles is the perpendicular bisector of the segment. Because the construction of the bisector is done without the knowledge of the segment's midpoint , the construction is used for determining as the intersection of the bisector and the line segment. This construction is in fact used when constructing a line perpendicular to a given line at a given point : drawing a circle whose center is such that it intersects the line in two points , and the perpendicular to be constructed is the one bisecting segment . Equations If are the position vectors of two points , then its midpoint is and vector is a normal vector of the perpendicular line segment bisector. Hence its vector equation is . Inserting and expanding the equation leads to the vector equation (V) With one gets the equation in coordinate form: (C) Or explicitly: (E), where , , and . Applications Perpendicular line segment bisectors were used solving various geometric problems: Construction of the center of a Thales' circle, Construction of the center of the Excircle of a triangle, Voronoi diagram boundaries consist of segments of such lines or planes. Perpendicular line segment bisectors in space The perpendicular bisector of a line segment is a plane, which meets the segment at its midpoint perpendicularly. Its vector equation is literally the same as in the plane case: (V) With one gets the equation in coordinate form: (C3) Property (D) (see above) is literally true in space, too: (D) The perpendicular bisector plane of a segment has for any point the property: . Angle bisector An angle bisector divides the angle into two angles with equal measures. An angle only has one bisector. Each point of an angle bisector is equidistant from the sides of the angle. The 'interior' or 'internal bisector' of an angle is the line, half-line, or line segment that divides an angle of less than 180° into two equal angles. The 'exterior' or 'external bisector' is the line that divides the supplementary angle (of 180° minus the original angle), formed by one side forming the original angle and the extension of the other side, into two equal angles. To bisect an angle with straightedge and compass, one draws a circle whose center is the vertex. The circle meets the angle at two points: one on each leg. Using each of these points as a center, draw two circles of the same size. The intersection of the circles (two points) determines a line that is the angle bisector. The proof of the correctness of this construction is fairly intuitive, relying on the symmetry of the problem. The trisection of an angle (dividing it into three equal parts) cannot be achieved with the compass and ruler alone (this was first proved by Pierre Wantzel). The internal and external bisectors of an angle are perpendicular. If the angle is formed by the two lines given algebraically as and then the internal and external bisectors are given by the two equations Triangle Concurrencies and collinearities The bisectors of two exterior angles and the bisector of the other interior angle are concurrent. Three intersection points, each of an external angle bisector with the opposite extended side, are collinear (fall on the same line as each other). Three intersection points, two of them between an interior angle bisector and the opposite side, and the third between the other exterior angle bisector and the opposite side extended, are collinear. Angle bisector theorem The angle bisector theorem is concerned with the relative lengths of the two segments that a triangle's side is divided into by a line that bisects the opposite angle. It equates their relative lengths to the relative lengths of the other two sides of the triangle. Lengths If the side lengths of a triangle are , the semiperimeter and A is the angle opposite side , then the length of the internal bisector of angle A is or in trigonometric terms, If the internal bisector of angle A in triangle ABC has length and if this bisector divides the side opposite A into segments of lengths m and n, then where b and c are the side lengths opposite vertices B and C; and the side opposite A is divided in the proportion b:c. If the internal bisectors of angles A, B, and C have lengths and , then No two non-congruent triangles share the same set of three internal angle bisector lengths. Integer triangles There exist integer triangles with a rational angle bisector. Quadrilateral The internal angle bisectors of a convex quadrilateral either form a cyclic quadrilateral (that is, the four intersection points of adjacent angle bisectors are concyclic), or they are concurrent. In the latter case the quadrilateral is a tangential quadrilateral. Rhombus Each diagonal of a rhombus bisects opposite angles. Ex-tangential quadrilateral The excenter of an ex-tangential quadrilateral lies at the intersection of six angle bisectors. These are the internal angle bisectors at two opposite vertex angles, the external angle bisectors (supplementary angle bisectors) at the other two vertex angles, and the external angle bisectors at the angles formed where the extensions of opposite sides intersect. Parabola The tangent to a parabola at any point bisects the angle between the line joining the point to the focus and the line from the point and perpendicular to the directrix. Bisectors of the sides of a polygon Triangle Medians Each of the three medians of a triangle is a line segment going through one vertex and the midpoint of the opposite side, so it bisects that side (though not in general perpendicularly). The three medians intersect each other at a point which is called the centroid of the triangle, which is its center of mass if it has uniform density; thus any line through a triangle's centroid and one of its vertices bisects the opposite side. The centroid is twice as close to the midpoint of any one side as it is to the opposite vertex. Perpendicular bisectors The interior perpendicular bisector of a side of a triangle is the segment, falling entirely on and inside the triangle, of the line that perpendicularly bisects that side. The three perpendicular bisectors of a triangle's three sides intersect at the circumcenter (the center of the circle through the three vertices). Thus any line through a triangle's circumcenter and perpendicular to a side bisects that side. In an acute triangle the circumcenter divides the interior perpendicular bisectors of the two shortest sides in equal proportions. In an obtuse triangle the two shortest sides' perpendicular bisectors (extended beyond their opposite triangle sides to the circumcenter) are divided by their respective intersecting triangle sides in equal proportions. For any triangle the interior perpendicular bisectors are given by and where the sides are and the area is Quadrilateral The two bimedians of a convex quadrilateral are the line segments that connect the midpoints of opposite sides, hence each bisecting two sides. The two bimedians and the line segment joining the midpoints of the diagonals are concurrent at a point called the "vertex centroid" and are all bisected by this point. The four "maltitudes" of a convex quadrilateral are the perpendiculars to a side through the midpoint of the opposite side, hence bisecting the latter side. If the quadrilateral is cyclic (inscribed in a circle), these maltitudes are concurrent at (all meet at) a common point called the "anticenter". Brahmagupta's theorem states that if a cyclic quadrilateral is orthodiagonal (that is, has perpendicular diagonals), then the perpendicular to a side from the point of intersection of the diagonals always bisects the opposite side. The perpendicular bisector construction forms a quadrilateral from the perpendicular bisectors of the sides of another quadrilateral. Area bisectors and perimeter bisectors Triangle There is an infinitude of lines that bisect the area of a triangle. Three of them are the medians of the triangle (which connect the sides' midpoints with the opposite vertices), and these are concurrent at the triangle's centroid; indeed, they are the only area bisectors that go through the centroid. Three other area bisectors are parallel to the triangle's sides; each of these intersects the other two sides so as to divide them into segments with the proportions . These six lines are concurrent three at a time: in addition to the three medians being concurrent, any one median is concurrent with two of the side-parallel area bisectors. The envelope of the infinitude of area bisectors is a deltoid (broadly defined as a figure with three vertices connected by curves that are concave to the exterior of the deltoid, making the interior points a non-convex set). The vertices of the deltoid are at the midpoints of the medians; all points inside the deltoid are on three different area bisectors, while all points outside it are on just one. The sides of the deltoid are arcs of hyperbolas that are asymptotic to the extended sides of the triangle. The ratio of the area of the envelope of area bisectors to the area of the triangle is invariant for all triangles, and equals i.e. 0.019860... or less than 2%. A cleaver of a triangle is a line segment that bisects the perimeter of the triangle and has one endpoint at the midpoint of one of the three sides. The three cleavers concur at (all pass through) the center of the Spieker circle, which is the incircle of the medial triangle. The cleavers are parallel to the angle bisectors. A splitter of a triangle is a line segment having one endpoint at one of the three vertices of the triangle and bisecting the perimeter. The three splitters concur at the Nagel point of the triangle. Any line through a triangle that splits both the triangle's area and its perimeter in half goes through the triangle's incenter (the center of its incircle). There are either one, two, or three of these for any given triangle. A line through the incenter bisects one of the area or perimeter if and only if it also bisects the other. Parallelogram Any line through the midpoint of a parallelogram bisects the area and the perimeter. Circle and ellipse All area bisectors and perimeter bisectors of a circle or other ellipse go through the center, and any chords through the center bisect the area and perimeter. In the case of a circle they are the diameters of the circle. Bisectors of diagonals Parallelogram The diagonals of a parallelogram bisect each other. Quadrilateral If a line segment connecting the diagonals of a quadrilateral bisects both diagonals, then this line segment (the Newton Line) is itself bisected by the vertex centroid. Volume bisectors A plane that divides two opposite edges of a tetrahedron in a given ratio also divides the volume of the tetrahedron in the same ratio. Thus any plane containing a bimedian (connector of opposite edges' midpoints) of a tetrahedron bisects the volume of the tetrahedron	Mathematics	Other	null
152561	https://en.wikipedia.org/wiki/Rutabaga	Rutabaga	Rutabaga (; North American English) or swede (English and some Commonwealth English) is a root vegetable, a form of Brassica napus (which also includes rapeseed). Other names include Swedish turnip, neep (Scots), and turnip (Scottish and Canadian English, Irish English and Manx English, as well as some dialects of English in Northern England). However, elsewhere, the name turnip usually refers to the related white turnip. The species Brassica napus originated as a hybrid between the cabbage (Brassica oleracea) and the turnip (Brassica rapa). Rutabaga roots are eaten as human food in various ways, and the leaves can be eaten as a leaf vegetable. The roots and tops are also used for livestock, fed directly in the winter or foraged in the field during the other seasons. Scotland, Northern and Western England, Wales, the Isle of Man, and Ireland had a tradition of carving the roots into Jack-o'-lanterns at Halloween. Etymology Rutabaga has many national and regional names. Rutabaga is the common North American term for the plant. This comes from the Swedish dialectal word , from 'root' + 'lump, bunch'. In the U.S., the plant is also known as Swedish turnip or yellow turnip. The term swede (from "Swedish turnip") is used in many Commonwealth Nations, including much of the United Kingdom, Australia, and New Zealand. The name turnip is also used in parts of Northern and Midland England, the West Country (particularly Cornwall), Ireland, the Isle of Man, and Canada. In Wales, according to region, it is variously known as , , or in Welsh, and as swede or turnip in English. In Scotland, it is known as turnip, tumshie (also used as a pejorative term for a foolish or stupid person), or neep (from Old English , Latin ). Some areas of south-east Scotland, such as Berwickshire and Roxburghshire, still use the term baigie, possibly a derivative of the Swedish dialectal word . The term turnip is also used for the white turnip (Brassica rapa ssp rapa). Some will also refer to both swede and (white) turnip as just turnip (this word is also derived from ). In north-east England, turnips and swedes are colloquially called snannies snadgers, snaggers (archaic) or narkies. Rutabaga is also known as moot in the Isle of Man and the Manx language word for turnip is . History The first known printed reference to the rutabaga comes from the Swiss botanist Gaspard Bauhin in 1620, where he notes that it was growing wild in Sweden. It is often considered to have originated in Scandinavia, Finland or Russia. According to the Natural Resources Institute Finland (now Luke), rutabaga or was most likely bred on more than one occasion in Northern Europe around the 16th century. Studies by its research institute have shown that was developed independently in Finland and Sweden from turnip and cabbage in connection with seed cultivation. There are contradictory accounts of how rutabaga arrived in England. Some sources say it arrived in England from Germany, while other accounts support Swedish origins. According to John Sinclair, the root vegetable arrived in England from Germany around 1750. Rutabaga arrived in Scotland by way of Sweden around 1781. An article in The Gardeners' Chronicle suggests that the rutabaga was introduced more widely to England in 1790. Introduction to North America came in the early 19th century with reports of rutabaga crops in Illinois as early as 1817. In 1835, a rutabaga fodder crop was recommended to New York farmers in the Genesee River valley. Rutabaga was once considered a food of last resort in both Germany and France due to its association with food shortages in World War I and World War II. Boiled stew with rutabaga and water as the only ingredients (Steckrübeneintopf) was a typical food in Germany during the famines and food shortages of World War I caused by the Allied blockade (the or Turnip Winter of 1916–17) and between 1945 and 1949. As a result, many older Germans had unhappy memories of this food. Botanical history Rutabaga has a complex taxonomic history. The earliest account comes from the Swiss botanist Gaspard Bauhin, who wrote about it in his 1620 Prodromus. Brassica napobrassica was first validly published by Carl Linnaeus in his 1753 work Species Plantarum as a variety of B. oleracea: B. oleracea var. napobrassica. It has since been moved to other taxa as a variety, subspecies, or elevated to species rank. In 1768, a Scottish botanist promoted Linnaeus' variety to species rank as Brassica napobrassica in The Gardeners Dictionary. Rutabaga has a chromosome number of 2n = 38. It originated from a cross between turnip (Brassica rapa) and Brassica oleracea. The resulting cross doubled its chromosomes, becoming an allopolyploid. This relationship was first published by Woo Jang-choon in 1935 and is known as the Triangle of U. Cuisine Europe Netherlands In the Netherlands, rutabaga is traditionally served boiled and mashed. Adding mashed potatoes (and, in some recipes, similarly mashed vegetables or fruits) makes 'mash pot', a dish often served alongside smoked sausage. Similar dishes are known in the Southern low countries, down to and including Brussels, as stoemp. Poland During the difficult days of World War II, rutabaga and rutabaga juice were an important part of the local diet, and were consumed in large quantities. Scandinavia Sweden and Norway In Sweden and Norway, rutabaga is cooked with potato and sometimes carrot, and mashed with butter and either stock or, occasionally, milk or cream, to create a puree called (Swedish, literally 'root mash') or (Norwegian). Onion is occasionally added. In Norway, is an obligatory accompaniment to many festive dishes, including , , and salted herring. In Sweden, is often eaten together with cured and boiled ham hock, accompanied by mustard. This classic Swedish dish is called . Finland Finns eat and cook rutabaga in a variety of ways. Rutabaga is the major ingredient in the popular Christmas dish lanttulaatikko (rutabaga casserole), one of the three main casseroles served during Finnish Christmas, alongside the potato and carrot casseroles. Uncooked and thinly julienned rutabaga is often served as a side dish salad in school and workplace lunches. Raisins or canned pineapple in light syrup are often added to the rutabaga salad. Sometimes, thinly sliced raw carrots are mixed with rutabaga. Finns use rutabaga in most dishes that call for a root vegetable. Many Finnish soup bases consist of potatoes, carrots, and rutabagas. Finnish cuisine also roasts, bakes, boils, and grills rutabagas. Oven-baked root vegetables are another home-cooking classic in Finland: rutabaga, carrots, beetroots, and potatoes are roasted in the oven with salt and oil. Karelian hot pot () is a popular slow-cooking stew with root vegetables and meat cooked for a long time in a Dutch oven. Finnish supermarkets sell alternative potato chips made from root vegetables, such as rutabagas, beetroots and carrots. Rutabagas are also an ingredient in (rutabaga-, a traditional Savonian and Karelian dish). United Kingdom England In England, swede is boiled with carrots and mashed or pureed with butter and ground pepper. The flavoured cooking water is often retained for soup or as an addition to gravy. Swede is also a component of the popular condiment Branston Pickle. The swede is also one of the four traditional ingredients of the pasty originating in Cornwall. Scotland In Scotland, separately boiled and mashed, swede (neeps) and potatoes are served as "neeps and " ( being the Scots word for potatoes), in a traditional Burns supper, together with the main course of haggis (the Scottish national dish). Neeps mashed with potatoes are called clapshot. Roughly equal quantities of neeps and tatties are boiled in salted water and mashed with butter. Seasoning can be augmented with black pepper. Onions are never used. Regionally, neeps are a common ingredient in soups and stews. Wales Swede is an essential vegetable component of the traditional Welsh lamb broth called cawl. A mash produced using just potato and swede is known as in the North-East of the country, as on the Llyn peninsula and as in other parts. Outside Europe Australia In Australia, swedes are used as a flavour enhancer in casseroles, stews, and soups. Canada In Canada, they are considered winter vegetables, as, along with similar vegetables, they can be kept in a cold area or cellar for several months. They are primarily used as a side dish. They are also used as filler in foods such as mincemeat and Christmas cake. In Newfoundland, it is served with Jiggs dinner. New Zealand In New Zealand, they are more commonly available in winter but can be easily purchased for much of the year. It is thought they best grow in Southland, where the winters are colder. They are usually served mashed with butter but are often added to other dishes like casseroles or bakes. United States In the US, rutabagas are not widely eaten but may be found as part of stews or casseroles, served mashed with carrots, or baked in a pasty. They are sometimes included in the New England boiled dinner. Phytochemistry Rutabaga and other cyanoglucoside-containing foods (including cassava, maize (corn), bamboo shoots, sweet potatoes, and lima beans) release cyanide, which is subsequently detoxified into thiocyanate. Thiocyanate inhibits thyroid iodide transport and, at high doses, competes with iodide in the organification process within thyroid tissue. Goitres may develop when there is a dietary imbalance of thiocyanate-containing food in excess of iodine consumption, and these compounds can contribute to hypothyroidism. Yet, there have been no reports of ill effects in humans from the consumption of glucosinolates from normal amounts of Brassica vegetables. Glucosinolate content in Brassica vegetables is around one percent of dry matter. These compounds also cause the bitter taste of rutabaga. As with watercress, mustard greens, turnip, broccoli, and horseradish, human perception of bitterness in rutabaga is governed by a gene affecting the TAS2R bitter receptor, which detects the glucosinolates in rutabaga. Sensitive individuals with the genotype PAV/PAV (supertasters) find rutabaga twice as bitter as insensitive subjects (AVI/AVI). The difference for the mixed type (PAV/AVI) is insignificant for rutabaga. As a result, sensitive individuals may find some rutabagas too bitter to eat. Other chemical compounds that contribute to flavour and odour include glucocheirolin, glucobrassicanapin, glucoberteroin, gluconapoleiferin, and glucoerysolin. Several phytoalexins that aid in defence against plant pathogens have also been isolated from the rutabaga, including three novel phytoalexins that were reported in 2004. Rutabaga contains significant amounts of vitamin C: 100 g contains 25 mg, 30% of the daily recommended dose. Other uses Livestock The roots and tops of "swedes" came into use as a forage crop in the early nineteenth century, used as winter feed for livestock. They may be fed directly (chopped or from a hopper), or animals may be allowed to forage the plants directly in the field. Halloween People in Northern England, West England, Ireland, and Scotland have long carved turnips and often use them as lanterns to ward off harmful spirits. In the Middle Ages, rowdy bands of children roamed the streets in masks carrying carved turnips known in Scotland as "tumshie heads". In modern times, turnips are often carved to look as sinister and threatening as possible and are put in the window or on the doorstep of a house on Halloween to ward off evil spirits. Since pumpkins became readily available in Europe in the 1980s, they have taken over this role to a large extent. In the Isle of Man, turnip lanterns are still carved at Hop-tu-Naa (Manx equivalent of Halloween), lit with a candle or electric torch, and carried from house to house by some children, with the accompanying Hop tu Naa song; hoping for money or treats of food. The smell of burning turnip is an evocative part of the event. Festivals A local farmers' market in the town of Ithaca, New York, organizes what it calls the International Rutabaga Curling Championship annually on the last day of the market season. The villages of Askov, Minnesota, and Cumberland, Wisconsin, both hold annual rutabaga festivals in August.	Biology and health sciences	Brassicales	null
152611	https://en.wikipedia.org/wiki/Cellular%20differentiation	Cellular differentiation	Cellular differentiation is the process in which a stem cell changes from one type to a differentiated one. Usually, the cell changes to a more specialized type. Differentiation happens multiple times during the development of a multicellular organism as it changes from a simple zygote to a complex system of tissues and cell types. Differentiation continues in adulthood as adult stem cells divide and create fully differentiated daughter cells during tissue repair and during normal cell turnover. Some differentiation occurs in response to antigen exposure. Differentiation dramatically changes a cell's size, shape, membrane potential, metabolic activity, and responsiveness to signals. These changes are largely due to highly controlled modifications in gene expression and are the study of epigenetics. With a few exceptions, cellular differentiation almost never involves a change in the DNA sequence itself. Metabolic composition, however, gets dramatically altered where stem cells are characterized by abundant metabolites with highly unsaturated structures whose levels decrease upon differentiation. Thus, different cells can have very different physical characteristics despite having the same genome. A specialized type of differentiation, known as terminal differentiation, is of importance in some tissues, including vertebrate nervous system, striated muscle, epidermis and gut. During terminal differentiation, a precursor cell formerly capable of cell division permanently leaves the cell cycle, dismantles the cell cycle machinery and often expresses a range of genes characteristic of the cell's final function (e.g. myosin and actin for a muscle cell). Differentiation may continue to occur after terminal differentiation if the capacity and functions of the cell undergo further changes. Among dividing cells, there are multiple levels of cell potency, which is the cell's ability to differentiate into other cell types. A greater potency indicates a larger number of cell types that can be derived. A cell that can differentiate into all cell types, including the placental tissue, is known as totipotent. In mammals, only the zygote and subsequent blastomeres are totipotent, while in plants, many differentiated cells can become totipotent with simple laboratory techniques. A cell that can differentiate into all cell types of the adult organism is known as pluripotent. Such cells are called meristematic cells in higher plants and embryonic stem cells in animals, though some groups report the presence of adult pluripotent cells. Virally induced expression of four transcription factors Oct4, Sox2, , and Klf4 (Yamanaka factors) is sufficient to create pluripotent (iPS) cells from adult fibroblasts. A multipotent cell is one that can differentiate into multiple different, but closely related cell types. Oligopotent cells are more restricted than multipotent, but can still differentiate into a few closely related cell types. Finally, unipotent cells can differentiate into only one cell type, but are capable of self-renewal. In cytopathology, the level of cellular differentiation is used as a measure of cancer progression. "Grade" is a marker of how differentiated a cell in a tumor is. Mammalian cell types Three basic categories of cells make up the mammalian body: germ cells, somatic cells, and stem cells. Each of the approximately 37.2 trillion (3.72x1013) cells in an adult human has its own copy or copies of the genome except certain cell types, such as red blood cells, that lack nuclei in their fully differentiated state. Most cells are diploid; they have two copies of each chromosome. Such cells, called somatic cells, make up most of the human body, such as skin and muscle cells. Cells differentiate to specialize for different functions. Germ line cells are any line of cells that give rise to gametes—eggs and sperm—and thus are continuous through the generations. Stem cells, on the other hand, have the ability to divide for indefinite periods and to give rise to specialized cells. They are best described in the context of normal human development. Development begins when a sperm fertilizes an egg and creates a single cell that has the potential to form an entire organism. In the first hours after fertilization, this cell divides into identical cells. In humans, approximately four days after fertilization and after several cycles of cell division, these cells begin to specialize, forming a hollow sphere of cells, called a blastocyst. The blastocyst has an outer layer of cells, and inside this hollow sphere, there is a cluster of cells called the inner cell mass. The cells of the inner cell mass go on to form virtually all of the tissues of the human body. Although the cells of the inner cell mass can form virtually every type of cell found in the human body, they cannot form an organism. These cells are referred to as pluripotent. Pluripotent stem cells undergo further specialization into multipotent progenitor cells that then give rise to functional cells. Examples of stem and progenitor cells include: Radial glial cells (embryonic neural stem cells) that give rise to excitatory neurons in the fetal brain through the process of neurogenesis. Hematopoietic stem cells (adult stem cells) from the bone marrow that give rise to red blood cells, white blood cells, and platelets. Mesenchymal stem cells (adult stem cells) from the bone marrow that give rise to stromal cells, fat cells, and types of bone cells Epithelial stem cells (progenitor cells) that give rise to the various types of skin cells Muscle satellite cells (progenitor cells) that contribute to differentiated muscle tissue. A pathway that is guided by the cell adhesion molecules consisting of four amino acids, arginine, glycine, asparagine, and serine, is created as the cellular blastomere differentiates from the single-layered blastula to the three primary layers of germ cells in mammals, namely the ectoderm, mesoderm and endoderm (listed from most distal (exterior) to proximal (interior)). The ectoderm ends up forming the skin and the nervous system, the mesoderm forms the bones and muscular tissue, and the endoderm forms the internal organ tissues. Dedifferentiation Dedifferentiation, or integration, is a cellular process seen in the more basal life forms in animals, such as worms and amphibians where a differentiated cell reverts to an earlier developmental stageusually as part of a regenerative process. Dedifferentiation also occurs in plant cells. And, in cell culture in the laboratory, cells can change shape or may lose specific properties such as protein expressionwhich processes are also termed dedifferentiation. Some hypothesize that dedifferentiation is an aberration that likely results in cancers, but others explain it as a natural part of the immune response that was lost to humans at some point of evolution. A newly discovered molecule dubbed reversine, a purine analog, has proven to induce dedifferentiation in myotubes. These manifestly dedifferentiated cellsnow performing essentially as stem cellscould then redifferentiate into osteoblasts and adipocytes. Mechanisms Each specialized cell type in an organism expresses a subset of all the genes that constitute the genome of that species. Each cell type is defined by its particular pattern of regulated gene expression. Cell differentiation is thus a transition of a cell from one cell type to another and it involves a switch from one pattern of gene expression to another. Cellular differentiation during development can be understood as the result of a gene regulatory network. A regulatory gene and its cis-regulatory modules are nodes in a gene regulatory network; they receive input and create output elsewhere in the network. The systems biology approach to developmental biology emphasizes the importance of investigating how developmental mechanisms interact to produce predictable patterns (morphogenesis). However, an alternative view has been proposed recently. Based on stochastic gene expression, cellular differentiation is the result of a Darwinian selective process occurring among cells. In this frame, protein and gene networks are the result of cellular processes and not their cause. While evolutionarily conserved molecular processes are involved in the cellular mechanisms underlying these switches, in animal species these are very different from the well-characterized gene regulatory mechanisms of bacteria, and even from those of the animals' closest unicellular relatives. Specifically, cell differentiation in animals is highly dependent on biomolecular condensates of regulatory proteins and enhancer DNA sequences. Cellular differentiation is often controlled by cell signaling. Many of the signal molecules that convey information from cell to cell during the control of cellular differentiation are called growth factors. Although the details of specific signal transduction pathways vary, these pathways often share the following general steps. A ligand produced by one cell binds to a receptor in the extracellular region of another cell, inducing a conformational change in the receptor. The shape of the cytoplasmic domain of the receptor changes, and the receptor acquires enzymatic activity. The receptor then catalyzes reactions that phosphorylate other proteins, activating them. A cascade of phosphorylation reactions eventually activates a dormant transcription factor or cytoskeletal protein, thus contributing to the differentiation process in the target cell. Cells and tissues can vary in competence, their ability to respond to external signals. Signal induction refers to cascades of signaling events, during which a cell or tissue signals to another cell or tissue to influence its developmental fate. Yamamoto and Jeffery investigated the role of the lens in eye formation in cave- and surface-dwelling fish, a striking example of induction. Through reciprocal transplants, Yamamoto and Jeffery found that the lens vesicle of surface fish can induce other parts of the eye to develop in cave- and surface-dwelling fish, while the lens vesicle of the cave-dwelling fish cannot. Other important mechanisms fall under the category of asymmetric cell divisions, divisions that give rise to daughter cells with distinct developmental fates. Asymmetric cell divisions can occur because of asymmetrically expressed maternal cytoplasmic determinants or because of signaling. In the former mechanism, distinct daughter cells are created during cytokinesis because of an uneven distribution of regulatory molecules in the parent cell; the distinct cytoplasm that each daughter cell inherits results in a distinct pattern of differentiation for each daughter cell. A well-studied example of pattern formation by asymmetric divisions is body axis patterning in Drosophila. RNA molecules are an important type of intracellular differentiation control signal. The molecular and genetic basis of asymmetric cell divisions has also been studied in green algae of the genus Volvox, a model system for studying how unicellular organisms can evolve into multicellular organisms. In Volvox carteri, the 16 cells in the anterior hemisphere of a 32-cell embryo divide asymmetrically, each producing one large and one small daughter cell. The size of the cell at the end of all cell divisions determines whether it becomes a specialized germ or somatic cell. Epigenetic control Since each cell, regardless of cell type, possesses the same genome, determination of cell type must occur at the level of gene expression. While the regulation of gene expression can occur through cis- and trans-regulatory elements including a gene's promoter and enhancers, the problem arises as to how this expression pattern is maintained over numerous generations of cell division. As it turns out, epigenetic processes play a crucial role in regulating the decision to adopt a stem, progenitor, or mature cell fate This section will focus primarily on mammalian stem cells. In systems biology and mathematical modeling of gene regulatory networks, cell-fate determination is predicted to exhibit certain dynamics, such as attractor-convergence (the attractor can be an equilibrium point, limit cycle or strange attractor) or oscillatory. Importance of epigenetic control The first question that can be asked is the extent and complexity of the role of epigenetic processes in the determination of cell fate. A clear answer to this question can be seen in the 2011 paper by Lister R, et al. on aberrant epigenomic programming in human induced pluripotent stem cells. As induced pluripotent stem cells (iPSCs) are thought to mimic embryonic stem cells in their pluripotent properties, few epigenetic differences should exist between them. To test this prediction, the authors conducted whole-genome profiling of DNA methylation patterns in several human embryonic stem cell (ESC), iPSC, and progenitor cell lines. Female adipose cells, lung fibroblasts, and foreskin fibroblasts were reprogrammed into induced pluripotent state with the OCT4, SOX2, KLF4, and MYC genes. Patterns of DNA methylation in ESCs, iPSCs, somatic cells were compared. Lister R, et al. observed significant resemblance in methylation levels between embryonic and induced pluripotent cells. Around 80% of CG dinucleotides in ESCs and iPSCs were methylated, the same was true of only 60% of CG dinucleotides in somatic cells. In addition, somatic cells possessed minimal levels of cytosine methylation in non-CG dinucleotides, while induced pluripotent cells possessed similar levels of methylation as embryonic stem cells, between 0.5 and 1.5%. Thus, consistent with their respective transcriptional activities, DNA methylation patterns, at least on the genomic level, are similar between ESCs and iPSCs. However, upon examining methylation patterns more closely, the authors discovered 1175 regions of differential CG dinucleotide methylation between at least one ES or iPS cell line. By comparing these regions of differential methylation with regions of cytosine methylation in the original somatic cells, 44-49% of differentially methylated regions reflected methylation patterns of the respective progenitor somatic cells, while 51-56% of these regions were dissimilar to both the progenitor and embryonic cell lines. In vitro-induced differentiation of iPSC lines saw transmission of 88% and 46% of hyper and hypo-methylated differentially methylated regions, respectively. Two conclusions are readily apparent from this study. First, epigenetic processes are heavily involved in cell fate determination, as seen from the similar levels of cytosine methylation between induced pluripotent and embryonic stem cells, consistent with their respective patterns of transcription. Second, the mechanisms of reprogramming (and by extension, differentiation) are very complex and cannot be easily duplicated, as seen by the significant number of differentially methylated regions between ES and iPS cell lines. Now that these two points have been established, we can examine some of the epigenetic mechanisms that are thought to regulate cellular differentiation. Mechanisms of epigenetic regulation Pioneer factors (Oct4, Sox2, Nanog) Three transcription factors, OCT4, SOX2, and NANOG – the first two of which are used in induced pluripotent stem cell (iPSC) reprogramming, along with Klf4 and c-Myc – are highly expressed in undifferentiated embryonic stem cells and are necessary for the maintenance of their pluripotency. It is thought that they achieve this through alterations in chromatin structure, such as histone modification and DNA methylation, to restrict or permit the transcription of target genes. While highly expressed, their levels require a precise balance to maintain pluripotency, perturbation of which will promote differentiation towards different lineages based on how the gene expression levels change. Differential regulation of Oct-4 and SOX2 levels have been shown to precede germ layer fate selection. Increased levels of Oct4 and decreased levels of Sox2 promote a mesendodermal fate, with Oct4 actively suppressing genes associated with a neural ectodermal fate. Similarly, increased levels of Sox2 and decreased levels of Oct4 promote differentiation towards a neural ectodermal fate, with Sox2 inhibiting differentiation towards a mesendodermal fate. Regardless of the lineage cells differentiate down, suppression of NANOG has been identified as a necessary prerequisite for differentiation. Polycomb repressive complex (PRC2) In the realm of gene silencing, Polycomb repressive complex 2, one of two classes of the Polycomb group (PcG) family of proteins, catalyzes the di- and tri-methylation of histone H3 lysine 27 (H3K27me2/me3). By binding to the H3K27me2/3-tagged nucleosome, PRC1 (also a complex of PcG family proteins) catalyzes the mono-ubiquitinylation of histone H2A at lysine 119 (H2AK119Ub1), blocking RNA polymerase II activity and resulting in transcriptional suppression. PcG knockout ES cells do not differentiate efficiently into the three germ layers, and deletion of the PRC1 and PRC2 genes leads to increased expression of lineage-affiliated genes and unscheduled differentiation. Presumably, PcG complexes are responsible for transcriptionally repressing differentiation and development-promoting genes. Trithorax group proteins (TrxG) Alternately, upon receiving differentiation signals, PcG proteins are recruited to promoters of pluripotency transcription factors. PcG-deficient ES cells can begin differentiation but cannot maintain the differentiated phenotype. Simultaneously, differentiation and development-promoting genes are activated by Trithorax group (TrxG) chromatin regulators and lose their repression. TrxG proteins are recruited at regions of high transcriptional activity, where they catalyze the trimethylation of histone H3 lysine 4 (H3K4me3) and promote gene activation through histone acetylation. PcG and TrxG complexes engage in direct competition and are thought to be functionally antagonistic, creating at differentiation and development-promoting loci what is termed a "bivalent domain" and rendering these genes sensitive to rapid induction or repression. DNA methylation Regulation of gene expression is further achieved through DNA methylation, in which the DNA methyltransferase-mediated methylation of cytosine residues in CpG dinucleotides maintains heritable repression by controlling DNA accessibility. The majority of CpG sites in embryonic stem cells are unmethylated and appear to be associated with H3K4me3-carrying nucleosomes. Upon differentiation, a small number of genes, including OCT4 and NANOG, are methylated and their promoters repressed to prevent their further expression. Consistently, DNA methylation-deficient embryonic stem cells rapidly enter apoptosis upon in vitro differentiation. Nucleosome positioning While the DNA sequence of most cells of an organism is the same, the binding patterns of transcription factors and the corresponding gene expression patterns are different. To a large extent, differences in transcription factor binding are determined by the chromatin accessibility of their binding sites through histone modification and/or pioneer factors. In particular, it is important to know whether a nucleosome is covering a given genomic binding site or not. This can be determined using a chromatin immunoprecipitation assay. Histone acetylation and methylation DNA-nucleosome interactions are characterized by two states: either tightly bound by nucleosomes and transcriptionally inactive, called heterochromatin, or loosely bound and usually, but not always, transcriptionally active, called euchromatin. The epigenetic processes of histone methylation and acetylation, and their inverses demethylation and deacetylation primarily account for these changes. The effects of acetylation and deacetylation are more predictable. An acetyl group is either added to or removed from the positively charged Lysine residues in histones by enzymes called histone acetyltransferases or histone deactylases, respectively. The acetyl group prevents Lysine's association with the negatively charged DNA backbone. Methylation is not as straightforward, as neither methylation nor demethylation consistently correlate with either gene activation or repression. However, certain methylations have been repeatedly shown to either activate or repress genes. The trimethylation of lysine 4 on histone 3 (H3K4Me3) is associated with gene activation, whereas trimethylation of lysine 27 on histone 3 represses genes In stem cells During differentiation, stem cells change their gene expression profiles. Recent studies have implicated a role for nucleosome positioning and histone modifications during this process. There are two components of this process: turning off the expression of embryonic stem cell (ESC) genes, and the activation of cell fate genes. Lysine specific demethylase 1 (KDM1A) is thought to prevent the use of enhancer regions of pluripotency genes, thereby inhibiting their transcription. It interacts with Mi-2/NuRD complex (nucleosome remodelling and histone deacetylase) complex, giving an instance where methylation and acetylation are not discrete and mutually exclusive, but intertwined processes. Role of signaling in epigenetic control A final question to ask concerns the role of cell signaling in influencing the epigenetic processes governing differentiation. Such a role should exist, as it would be reasonable to think that extrinsic signaling can lead to epigenetic remodeling, just as it can lead to changes in gene expression through the activation or repression of different transcription factors. Little direct data is available concerning the specific signals that influence the epigenome, and the majority of current knowledge about the subject consists of speculations on plausible candidate regulators of epigenetic remodeling. We will first discuss several major candidates thought to be involved in the induction and maintenance of both embryonic stem cells and their differentiated progeny, and then turn to one example of specific signaling pathways in which more direct evidence exists for its role in epigenetic change. The first major candidate is Wnt signaling pathway. The Wnt pathway is involved in all stages of differentiation, and the ligand Wnt3a can substitute for the overexpression of c-Myc in the generation of induced pluripotent stem cells. On the other hand, disruption of β-catenin, a component of the Wnt signaling pathway, leads to decreased proliferation of neural progenitors. Growth factors comprise the second major set of candidates of epigenetic regulators of cellular differentiation. These morphogens are crucial for development, and include bone morphogenetic proteins, transforming growth factors (TGFs), and fibroblast growth factors (FGFs). TGFs and FGFs have been shown to sustain expression of OCT4, SOX2, and NANOG by downstream signaling to Smad proteins. Depletion of growth factors promotes the differentiation of ESCs, while genes with bivalent chromatin can become either more restrictive or permissive in their transcription. Several other signaling pathways are also considered to be primary candidates. Cytokine leukemia inhibitory factors are associated with the maintenance of mouse ESCs in an undifferentiated state. This is achieved through its activation of the Jak-STAT3 pathway, which has been shown to be necessary and sufficient towards maintaining mouse ESC pluripotency. Retinoic acid can induce differentiation of human and mouse ESCs, and Notch signaling is involved in the proliferation and self-renewal of stem cells. Finally, Sonic hedgehog, in addition to its role as a morphogen, promotes embryonic stem cell differentiation and the self-renewal of somatic stem cells. The problem, of course, is that the candidacy of these signaling pathways was inferred primarily on the basis of their role in development and cellular differentiation. While epigenetic regulation is necessary for driving cellular differentiation, they are certainly not sufficient for this process. Direct modulation of gene expression through modification of transcription factors plays a key role that must be distinguished from heritable epigenetic changes that can persist even in the absence of the original environmental signals. Only a few examples of signaling pathways leading to epigenetic changes that alter cell fate currently exist, and we will focus on one of them. Expression of Shh (Sonic hedgehog) upregulates the production of BMI1, a component of the PcG complex that recognizes H3K27me3. This occurs in a Gli-dependent manner, as Gli1 and Gli2 are downstream effectors of the Hedgehog signaling pathway. In culture, Bmi1 mediates the Hedgehog pathway's ability to promote human mammary stem cell self-renewal. In both humans and mice, researchers showed Bmi1 to be highly expressed in proliferating immature cerebellar granule cell precursors. When Bmi1 was knocked out in mice, impaired cerebellar development resulted, leading to significant reductions in postnatal brain mass along with abnormalities in motor control and behavior. A separate study showed a significant decrease in neural stem cell proliferation along with increased astrocyte proliferation in Bmi null mice. An alternative model of cellular differentiation during embryogenesis is that positional information is based on mechanical signalling by the cytoskeleton using Embryonic differentiation waves. The mechanical signal is then epigenetically transduced via signal transduction systems (of which specific molecules such as Wnt are part) to result in differential gene expression. In summary, the role of signaling in the epigenetic control of cell fate in mammals is largely unknown, but distinct examples exist that indicate the likely existence of further such mechanisms. Effect of matrix elasticity In order to fulfill the purpose of regenerating a variety of tissues, adult stems are known to migrate from their niches, adhere to new extracellular matrices (ECM) and differentiate. The ductility of these microenvironments are unique to different tissue types. The ECM surrounding brain, muscle and bone tissues range from soft to stiff. The transduction of the stem cells into these cells types is not directed solely by chemokine cues and cell to cell signaling. The elasticity of the microenvironment can also affect the differentiation of mesenchymal stem cells (MSCs which originate in bone marrow.) When MSCs are placed on substrates of the same stiffness as brain, muscle and bone ECM, the MSCs take on properties of those respective cell types. Matrix sensing requires the cell to pull against the matrix at focal adhesions, which triggers a cellular mechano-transducer to generate a signal to be informed what force is needed to deform the matrix. To determine the key players in matrix-elasticity-driven lineage specification in MSCs, different matrix microenvironments were mimicked. From these experiments, it was concluded that focal adhesions of the MSCs were the cellular mechano-transducer sensing the differences of the matrix elasticity. The non-muscle myosin IIa-c isoforms generates the forces in the cell that lead to signaling of early commitment markers. Nonmuscle myosin IIa generates the least force increasing to non-muscle myosin IIc. There are also factors in the cell that inhibit non-muscle myosin II, such as blebbistatin. This makes the cell effectively blind to the surrounding matrix. Researchers have achieved some success in inducing stem cell-like properties in HEK 239 cells by providing a soft matrix without the use of diffusing factors. The stem-cell properties appear to be linked to tension in the cells' actin network. One identified mechanism for matrix-induced differentiation is tension-induced proteins, which remodel chromatin in response to mechanical stretch. The RhoA pathway is also implicated in this process. Evolutionary history A billion-years-old, likely holozoan, protist, Bicellum brasieri with two types of cells, shows that the evolution of differentiated multicellularity, possibly but not necessarily of animal lineages, occurred at least 1 billion years ago and possibly mainly in freshwater lakes rather than the ocean.	Biology and health sciences	Cell processes	Biology
152623	https://en.wikipedia.org/wiki/Radiology	Radiology	Radiology ( ) is the medical specialty that uses medical imaging to diagnose diseases and guide treatment within the bodies of humans and other animals. It began with radiography (which is why its name has a root referring to radiation), but today it includes all imaging modalities. This includes technologies that use no ionizing electromagnetic radiation, such as ultrasonography and magnetic resonance imaging), as well as others that do use radiation, such as computed tomography (CT), fluoroscopy, and nuclear medicine including positron emission tomography (PET). Interventional radiology is the performance of usually minimally invasive medical procedures with the guidance of imaging technologies such as those mentioned above. The modern practice of radiology involves a team of several different healthcare professionals. A radiologist, who is a medical doctor with specialized post-graduate training, interprets medical images, communicates these findings to other physicians through reports or verbal communication, and uses imaging to perform minimally invasive medical procedures The nurse is involved in the care of patients before and after imaging or procedures, including administration of medications, monitoring of vital signs and monitoring of sedated patients. The radiographer, also known as a "radiologic technologist" in some countries such as the United States and Canada, is a specially trained healthcare professional that uses sophisticated technology and positioning techniques to produce medical images for the radiologist to interpret. Depending on the individual's training and country of practice, the radiographer may specialize in one of the above-mentioned imaging modalities or have expanded roles in image reporting. Diagnostic imaging modalities Projection (plain) radiography Radiographs (originally called roentgenographs, named after the discoverer of X-rays, Wilhelm Conrad Röntgen) are produced by transmitting X-rays through a patient. The X-rays are projected through the body onto a detector; an image is formed based on which rays pass through (and are detected) versus those that are absorbed or scattered in the patient (and thus are not detected). Röntgen discovered X-rays on November 8, 1895, and received the first Nobel Prize in Physics for his discovery in 1901. In film-screen radiography, an X-ray tube generates a beam of X-rays, which is aimed at the patient. The X-rays that pass through the patient are filtered through a device called a grid or X-ray filter, to reduce scatter, and strike an undeveloped film, which is held tightly to a screen of light-emitting phosphors in a light-tight cassette. The film is then developed chemically and an image appears on the film. Film-screen radiography is being replaced by phosphor plate radiography but more recently by digital radiography (DR) and the EOS imaging. In the two latest systems, the X-rays strike sensors that converts the signals generated into digital information, which is transmitted and converted into an image displayed on a computer screen. In digital radiography the sensors shape a plate, but in the EOS system, which is a slot-scanning system, a linear sensor vertically scans the patient. Plain radiography was the only imaging modality available during the first 50 years of radiology. Due to its availability, speed, and lower costs compared to other modalities, radiography is often the first-line test of choice in radiologic diagnosis. Also despite the large amount of data in CT scans, MR scans and other digital-based imaging, there are many disease entities in which the classic diagnosis is obtained by plain radiographs. Examples include various types of arthritis and pneumonia, bone tumors (especially benign bone tumors), fractures, congenital skeletal anomalies, and certain kidney stones. Mammography and DXA are two applications of low energy projectional radiography, used for the evaluation for breast cancer and osteoporosis, respectively. Fluoroscopy Fluoroscopy and angiography are special applications of X-ray imaging, in which a fluorescent screen and image intensifier tube is connected to a closed-circuit television system. This allows real-time imaging of structures in motion or augmented with a radiocontrast agent. Radiocontrast agents are usually administered by swallowing or injecting into the body of the patient to delineate anatomy and functioning of the blood vessels, the genitourinary system, or the gastrointestinal tract (GI tract). Two radiocontrast agents are presently in common use. Barium sulfate (BaSO4) is given orally or rectally for evaluation of the GI tract. Iodine, in multiple proprietary forms, is given by oral, rectal, vaginal, intra-arterial or intravenous routes. These radiocontrast agents strongly absorb or scatter X-rays, and in conjunction with the real-time imaging, allow demonstration of dynamic processes, such as peristalsis in the digestive tract or blood flow in arteries and veins. Iodine contrast may also be concentrated in abnormal areas more or less than in normal tissues and make abnormalities (tumors, cysts, inflammation) more conspicuous. Additionally, in specific circumstances, air can be used as a contrast agent for the gastrointestinal system and carbon dioxide can be used as a contrast agent in the venous system; in these cases, the contrast agent attenuates the X-ray radiation less than the surrounding tissues. Computed tomography CT imaging uses X-rays in conjunction with computing algorithms to image the body. In CT, an X-ray tube opposite an X-ray detector (or detectors) in a ring-shaped apparatus rotate around a patient, producing a computer-generated cross-sectional image (tomogram). CT is acquired in the axial plane, with coronal and sagittal images produced by computer reconstruction. Radiocontrast agents are often used with CT for enhanced delineation of anatomy. Although radiographs provide higher spatial resolution, CT can detect more subtle variations in attenuation of X-rays (higher contrast resolution). CT exposes the patient to significantly more ionizing radiation than a radiograph. Spiral multidetector CT uses 16, 64, 254 or more detectors during continuous motion of the patient through the radiation beam to obtain fine detail images in a short exam time. With rapid administration of intravenous contrast during the CT scan, these fine detail images can be reconstructed into three-dimensional (3D) images of carotid, cerebral, coronary or other arteries. The introduction of computed tomography in the early 1970s revolutionized diagnostic radiology by providing front-line clinicians with detailed images of anatomic structures in three dimensions. CT scanning has become the test of choice in diagnosing some urgent and emergent conditions, such as cerebral hemorrhage, pulmonary embolism (clots in the arteries of the lungs), aortic dissection (tearing of the aortic wall), appendicitis, diverticulitis, and obstructing kidney stones. Before the development of CT imaging, risky and painful exploratory surgery was often the only way to obtain a definitive diagnosis of the cause of severe abdominal pain which could not be otherwise ascertained from external observation. Continuing improvements in CT technology, including faster scanning times and improved resolution, have dramatically increased the accuracy and usefulness of CT scanning, which may partially account for increased use in medical diagnosis. Ultrasound Medical ultrasonography uses ultrasound (high-frequency sound waves) to visualize soft tissue structures in the body in real time. No ionizing radiation is involved, but the quality of the images obtained using ultrasound is highly dependent on the skill of the person (ultrasonographer) performing the exam and the patient's body size. Examinations of larger, overweight patients may have a decrease in image quality as their subcutaneous fat absorbs more of the sound waves. This results in fewer sound waves penetrating to organs and reflecting back to the transducer, resulting in loss of information and a poorer quality image. Ultrasound is also limited by its inability to image through air pockets (lungs, bowel loops) or bone. Its use in medical imaging has developed mostly within the last 30 years. The first ultrasound images were static and two-dimensional (2D), but with modern ultrasonography, 3D reconstructions can be observed in real time, effectively becoming "4D". Because ultrasound imaging techniques do not employ ionizing radiation to generate images (unlike radiography, and CT scans), they are generally considered safer and are therefore more common in obstetrical imaging. The progression of pregnancies can be thoroughly evaluated with less concern about damage from the techniques employed, allowing early detection and diagnosis of many fetal anomalies. Growth can be assessed over time, important in patients with chronic disease or pregnancy-induced disease, and in multiple pregnancies (twins, triplets, etc.). Color-flow Doppler ultrasound measures the severity of peripheral vascular disease and is used by cardiologists for dynamic evaluation of the heart, heart valves and major vessels. Stenosis, for example, of the carotid arteries may be a warning sign for an impending stroke. A clot, embedded deep in one of the inner veins of the legs, can be found via ultrasound before it dislodges and travels to the lungs, resulting in a potentially fatal pulmonary embolism. Ultrasound is useful as a guide to performing biopsies to minimize damage to surrounding tissues and in drainages such as thoracentesis. Small, portable ultrasound devices now replace peritoneal lavage in trauma wards by non-invasively assessing for the presence of internal bleeding and any internal organ damage. Extensive internal bleeding or injury to the major organs may require surgery and repair. Magnetic resonance imaging MRI uses strong magnetic fields to align atomic nuclei (usually hydrogen protons) within body tissues, then uses a radio signal to disturb the axis of rotation of these nuclei and observes the radio frequency signal generated as the nuclei return to their baseline states. The radio signals are collected by small antennae, called coils, placed near the area of interest. An advantage of MRI is its ability to produce images in axial, coronal, sagittal and multiple oblique planes with equal ease. MRI scans give the best soft tissue contrast of all the imaging modalities. With advances in scanning speed and spatial resolution, and improvements in computer 3D algorithms and hardware, MRI has become an important tool in musculoskeletal radiology and neuroradiology. One disadvantage is the patient has to hold still for long periods of time in a noisy, cramped space while the imaging is performed. Claustrophobia (fear of closed spaces) severe enough to terminate the MRI exam is reported in up to 5% of patients. Recent improvements in magnet design including stronger magnetic fields (3 teslas), shortening exam times, wider, shorter magnet bores and more open magnet designs, have brought some relief for claustrophobic patients. However, for magnets with equivalent field strengths, there is often a trade-off between image quality and open design. MRI has great benefit in imaging the brain, spine, and musculoskeletal system. The use of MRI is currently contraindicated for patients with pacemakers, cochlear implants, some indwelling medication pumps, certain types of cerebral aneurysm clips, metal fragments in the eyes, some metallic hardware due to the powerful magnetic fields, and strong fluctuating radio signals to which the body is exposed. Areas of potential advancement include functional imaging, cardiovascular MRI, and MRI-guided therapy. Nuclear medicine Nuclear medicine imaging involves the administration into the patient of radiopharmaceuticals consisting of substances with affinity for certain body tissues labeled with radioactive tracer. The most commonly used tracers are technetium-99m, iodine-123, iodine-131, gallium-67, indium-111, thallium-201 and fludeoxyglucose (18F) (18F-FDG). The heart, lungs, thyroid, liver, brain, gallbladder, and bones are commonly evaluated for particular conditions using these techniques. While anatomical detail is limited in these studies, nuclear medicine is useful in displaying physiological function. The excretory function of the kidneys, iodine-concentrating ability of the thyroid, blood flow to heart muscle, etc. can be measured. The principal imaging devices are the gamma camera and the PET Scanner, which detect the radiation emitted by the tracer in the body and display it as an image. With computer processing, the information can be displayed as axial, coronal and sagittal images (single-photon emission computed tomography - SPECT or Positron-emission tomography - PET). In the most modern devices, nuclear medicine images can be fused with a CT scan taken quasisimultaneously, so the physiological information can be overlaid or coregistered with the anatomical structures to improve diagnostic accuracy. Positron emission tomography (PET) scanning deals with positrons instead of gamma rays detected by gamma cameras. The positrons annihilate to produce two opposite traveling gamma rays to be detected coincidentally, thus improving resolution. In PET scanning, a radioactive, biologically active substance, most often 18F-FDG, is injected into a patient and the radiation emitted by the patient is detected to produce multiplanar images of the body. Metabolically more active tissues, such as cancer, concentrate the active substance more than normal tissues. PET images can be combined (or "fused") with anatomic (CT) imaging, to more accurately localize PET findings and thereby improve diagnostic accuracy. The fusion technology has gone further to combine PET and MRI similar to PET and CT. PET/MRI fusion, largely practiced in academic and research settings, could potentially play a crucial role in fine detail of brain imaging, breast cancer screening, and small joint imaging of the foot. The technology recently blossomed after passing the technical hurdle of altered positron movement in strong magnetic field thus affecting the resolution of PET images and attenuation correction. Interventional radiology Interventional radiology (IR or sometimes VIR for vascular and interventional radiology) is a subspecialty of radiology in which minimally invasive procedures are performed using image guidance. Some of these procedures are done for purely diagnostic purposes (e.g., angiogram), while others are done for treatment purposes (e.g., angioplasty). The basic concept behind interventional radiology is to diagnose or treat pathologies, with the most minimally invasive technique possible. Minimally invasive procedures are currently performed more than ever before. These procedures are often performed with the patient fully awake, with little or no sedation required. Interventional radiologists and interventional radiographers diagnose and treat several disorders, including peripheral vascular disease, renal artery stenosis, inferior vena cava filter placement, gastrostomy tube placements, biliary stents and hepatic interventions. Radiographic images, fluoroscopy, and ultrasound modalities are used for guidance, and the primary instruments used during the procedure are specialized needles and catheters. The images provide maps that allow the clinician to guide these instruments through the body to the areas containing disease. By minimizing the physical trauma to the patient, peripheral interventions can reduce infection rates and recovery times, as well as hospital stays. To be a trained interventionalist in the United States, an individual completes a five-year residency in radiology and a one- or two-year fellowship in IR. Analysis of images Plain, or general, radiography The basic technique is optical density evaluation (i.e. histogram analysis). It is then described that a region has a different optical density, e.g. a cancer metastasis to bone can cause radiolucency. The development of this is the digital radiological subtraction. It consists in overlapping two radiographs of the same examined region and subtracting the optical densities Comparison of changes in dental and bone radiographic densities in the presence of different soft-tissue simulators using pixel intensity and digital subtraction analyses. The resultant image only contains the time-dependent differences between the two examined radiographs. The advantage of this technique is the precise determination of the dynamics of density changes and the place of their occurrence. However, beforehand the geometrical adjustment and general alignment of optical density should be done Noise in subtraction images made from pairs of intraoral radiographs: a comparison between four methods of geometric alignment. Another possibility of radiographic image analysis is to study second order features, e.g. digital texture analysis Basic research Textural entropy as a potential feature for quantitative assessment of jaw bone healing process Comparative Analysis of Three Bone Substitute Materials Based on Co-Occurrence Matrix or fractal dimension Using fractal dimension to evaluate alveolar bone defects treated with various bone substitute materials. On this basis, it is possible to assess the places where bio-materials are implanted into the bone for the purpose of guided bone regeneration. They take an intact bone image sample (region of interest, ROI, reference site) and a sample of the implantation site (second ROI, test site) can be assessed numerically/objectively to what extent the implantation site imitates a healthy bone and how advanced is the process of bone regeneration Fast-Versus Slow-Resorbable Calcium Phosphate Bone Substitute Materials—Texture Analysis after 12 Months of Observation New Oral Surgery Materials for Bone Reconstruction—A Comparison of Five Bone Substitute Materials for Dentoalveolar Augmentation. It is also possible to check whether the bone healing process is influenced by some systemic factors Influence of General Mineral Condition on Collagen-Guided Alveolar Crest Augmentation. Teleradiology Teleradiology is the transmission of radiographic images from one location to another for interpretation by an appropriately trained professional, usually a radiologist or reporting radiographer. It is most often used to allow rapid interpretation of emergency room, ICU and other emergent examinations after hours of usual operation, at night and on weekends. In these cases, the images can be sent across time zones (e.g. to Spain, Australia, India) with the receiving Clinician working his normal daylight hours. However, at present, large private teleradiology companies in the U.S. currently provide most after-hours coverage employing night-working radiologists in the U.S. Teleradiology can also be used to obtain consultation with an expert or subspecialist about a complicated or puzzling case. In the U.S., many hospitals outsource their radiology departments to radiologists in India due to the lowered cost and availability of high speed internet access. Teleradiology requires a sending station, a high-speed internet connection, and a high-quality receiving station. At the transmission station, plain radiographs are passed through a digitizing machine before transmission, while CT, MRI, ultrasound and nuclear medicine scans can be sent directly, as they are already digital data. The computer at the receiving end will need to have a high-quality display screen that has been tested and cleared for clinical purposes. Reports are then transmitted to the requesting clinician. The major advantage of teleradiology is the ability to use different time zones to provide real-time emergency radiology services around-the-clock. The disadvantages include higher costs, limited contact between the referrer and the reporting Clinician, and the inability to cover for procedures requiring an onsite reporting Clinician. Laws and regulations concerning the use of teleradiology vary among the states, with some requiring a license to practice medicine in the state sending the radiologic exam. In the U.S., some states require the teleradiology report to be preliminary with the official report issued by a hospital staff radiologist. Lastly, a benefit of teleradiology is that it might be automated with modern machine learning techniques. Patient interaction Some radiologists, like teleradiologists, have no interaction with patients. Other radiologists, like interventional radiologists, primarily interact with patients and spend less time analyzing images. Diagnostic radiologists tend to spend the majority of their time analyzing images and a minority of their time interacting with patients. Compared to the healthcare provider who sends the patient to have images interpreted by a diagnostic radiologist, the radiologist usually does not know as much about the patient's clinical status or have as much influence on what action should be taken based on the images. Thus, the diagnostic radiologist reports image findings directly to that healthcare provider and often provides recommendations, who then takes the appropriate next steps for recommendations about medical management. Because radiologists undergo training regarding risks associated with different types of imaging tests and image-guided procedures, radiologists are the healthcare providers who generally educate patients about those risks to enable informed consent, not the healthcare provider requesting the test or procedure. Professional training United States Radiology is a field in medicine that has expanded rapidly after 2000 due to advances in computer technology, which is closely linked to modern imaging techniques. Applying for residency positions in radiology has become highly competitive. Applicants are often near the top of their medical school classes, with high USMLE (board) examination scores. Diagnostic radiologists must complete prerequisite undergraduate education, four years of medical school to earn a medical degree (D.O. or M.D.), one year of internship, and four years of residency training. After residency, most radiologists pursue one or two years of additional specialty fellowship training. The American Board of Radiology (ABR) administers professional certification in Diagnostic Radiology, Radiation Oncology and Medical Physics as well as subspecialty certification in neuroradiology, nuclear radiology, pediatric radiology and vascular and interventional radiology. "Board Certification" in diagnostic radiology requires successful completion of two examinations. The Core Exam is given after 36 months of residency. Although previously taken in Chicago or Tucson, Arizona, beginning in February 2021, the computer test transitioned permanently to a remote format. It encompasses 18 categories. A passing score is 350 or above. A fail on one to five categories was previously a Conditioned exam, however beginning in June 2021, the conditioned category will no longer exist and the test will be graded as a whole. The Certification Exam, can be taken 15 months after completion of the Radiology residency. This computer-based examination consists of five modules and graded pass-fail. It is given twice a year in Chicago and Tucson. Recertification examinations are taken every 10 years, with additional required continuing medical education as outlined in the Maintenance of Certification document. Certification may also be obtained from the American Osteopathic Board of Radiology (AOBR) and the American Board of Physician Specialties. Following completion of residency training, radiologists may either begin practicing as a general diagnostic radiologist or enter into subspecialty training programs known as fellowships. Examples of subspeciality training in radiology include abdominal imaging, thoracic imaging, cross-sectional/ultrasound, MRI, musculoskeletal imaging, interventional radiology, neuroradiology, interventional neuroradiology, paediatric radiology, nuclear medicine, emergency radiology, breast imaging and women's imaging. Fellowship training programs in radiology are usually one or two years in length. Some medical schools in the US have started to incorporate a basic radiology introduction into their core MD training. New York Medical College, the Wayne State University School of Medicine, Weill Cornell Medicine, the Uniformed Services University, and the University of South Carolina School of Medicine offer an introduction to radiology during their respective MD programs. Campbell University School of Osteopathic Medicine also integrates imaging material into their curriculum early in the first year. Radiographic exams are usually performed by radiographers. Qualifications for radiographers vary by country, but many radiographers now are required to hold a degree. Veterinary radiologists are veterinarians who specialize in the use of X-rays, ultrasound, MRI and nuclear medicine for diagnostic imaging or treatment of disease in animals. They are certified in either diagnostic radiology or radiation oncology by the American College of Veterinary Radiology. United Kingdom Radiology is an extremely competitive speciality in the UK, attracting applicants from a broad range of backgrounds. Applicants are welcomed directly from the Foundation Programme, as well as those who have completed higher training. Recruitment and selection into training post in clinical radiology posts in England, Scotland and Wales is done by an annual nationally coordinated process lasting from November to March. In this process, all applicants are required to pass a Specialty Recruitment Assessment (SRA) test. Those with a test score above a certain threshold are offered a single interview at the London and the South East Recruitment Office. At a later stage, applicants declare what programs they prefer, but may in some cases be placed in a neighbouring region. The training programme lasts for a total of five years. During this time, doctors rotate into different subspecialities, such as paediatrics, musculoskeletal or neuroradiology, and breast imaging. During the first year of training, radiology trainees are expected to pass the first part of the Fellowship of the Royal College of Radiologists (FRCR) exam. This comprises a medical physics and anatomy examination. Following completion of their part 1 exam, they are then required to pass six written exams (part 2A), which cover all the subspecialities. Successful completion of these allows them to complete the FRCR by completing part 2B, which includes rapid reporting, and a long case discussion. After achieving a certificate of completion of training (CCT), many fellowship posts exist in specialities such as neurointervention and vascular intervention, which would allow the doctor to work as an Interventional radiologist. In some cases, the CCT date can be deferred by a year to include these fellowship programmes. UK radiology registrars are represented by the Society of Radiologists in Training (SRT), which was founded in 1993 under the auspices of the Royal College of Radiologists. The society is a nonprofit organisation, run by radiology registrars specifically to promote radiology training and education in the UK. Annual meetings are held by which trainees across the country are encouraged to attend. Currently, a shortage of radiologists in the UK has created opportunities in all specialities, and with the increased reliance on imaging, demand is expected to increase in the future. Radiographers, and less frequently Nurses, are often trained to undertake many of these opportunities in order to help meet demand. Radiographers often may control a "list" of a particular set of procedures after being approved locally and signed off by a consultant radiologist. Similarly, radiographers may simply operate a list for a radiologist or other physician on their behalf. Most often if a radiographer operates a list autonomously then they are acting as the operator and practitioner under the Ionising Radiation (Medical Exposures) Regulations 2000. Radiographers are represented by a variety of bodies; most often this is the Society and College of Radiographers. Collaboration with nurses is also common, where a list may be jointly organised between the nurse and radiographer. Germany After obtaining medical licensure, German radiologists complete a five-year residency, culminating with a board examination (known as Facharztprüfung). Italy Italian radiologists complete a four-year residency program, after completing the six-year MD program. The Netherlands Dutch radiologists complete a five-year residency program, after completing the six-year MD program. India In India, one must obtain a bachelor's degree which requires 4.5 years of training, along with 1 year internship, followed by NEET PG examination which is one of the hardest examinations in India. Previous rank data shows only top rankers take radiology which means if the score is less, one might get accepted into other branches, but not radiology. The radiology program is a post graduate 3-year program (MD/DNB Radiology) or a 2-year diploma (DMRD). Singapore Radiologists in Singapore complete a five-year undergraduate MD program, followed by a one-year internship, and then a five-year residency program. Some radiologists may elect to complete a one or two-year fellowship for further sub-specialization in fields such as interventional radiology. Slovenia After finishing a six-year study of medicine and passing the emergency medicine internship, MDs can apply for radiology residency. Radiology is a five-year post-graduate program that involves all fields of radiology with a final board exam. France To become a radiologist, after having validated the common core of medical studies, one must obtain a DES (Specialized Studies Diploma) in radiology and medical imaging (specialized studies in 5 years), or a DES in advanced interventional radiology (specialized studies in 6 years). At the end of his DES, once validated, the future doctor will have to defend his “practice thesis” in order to validate his DE (State Diploma) as a doctor of medicine (common to all doctors of medicine therefore) and to be able to practice in France. Specialty training for interventional radiology Training for interventional radiology occurs in the residency portion of medical education, and has gone through developments. In 2000, the Society of Interventional Radiology (SIR) created a program named "Clinical Pathway in IR", which modified the "Holman Pathway" that was already accepted by the American Board of Radiology to including training in IR; this was accepted by ABR but was not widely adopted. In 2005 SIR proposed and ABR accepted another pathway called "DIRECT (Diagnostic and Interventional Radiology Enhanced Clinical Training) Pathway" to help trainees coming from other specialities learn IR; this too was not widely adopted. In 2006 SIR proposed a pathway resulting in certification in IR as a speciality; this was eventually accepted by the ABR in 2007 and was presented to the American Board of Medical Specialities (ABMS) in 2009, which rejected it because it did not include enough diagnostic radiology (DR) training. The proposal was reworked, at the same time that overall DR training was being revamped, and a new proposal that would lead to a dual DR/IR specialization was presented to the ABMS and was accepted in 2012 and eventually was implemented in 2014. By 2016 the field had determined that the old IR fellowships would be terminated by 2020. A handful of programs have offered interventional radiology fellowships that focus on training in the treatment of children. In Europe the field followed its own pathway; for example in Germany the parallel interventional society began to break free of the DR society in 2008. In the UK, interventional radiology was approved as a sub-specialty of clinical radiology in 2010. While many countries have an interventional radiology society, there is also the European-wide Cardiovascular and Interventional Radiological Society of Europe, whose aim is to support teaching, science, research and clinical practice in the field by hosting meetings, educational workshops and promoting patient safety initiatives. Furthermore, the Society provides an examination, the European Board of Interventional Radiology (EBIR), which is a highly valuable qualification in interventional radiology based on the European Curriculum and Syllabus for IR.	Biology and health sciences	Fields of medicine	null
152630	https://en.wikipedia.org/wiki/Speleology	Speleology	Speleology () is the scientific study of caves and other karst features, as well as their composition, structure, physical properties, history, ecology, and the processes by which they form (speleogenesis) and change over time (speleomorphology). The term speleology is also sometimes applied to the recreational activity of exploring caves, but this is more properly known as caving, potholing (British English), or spelunking (United States and Canadian English). Speleology and caving are often connected, as the physical skills required for in situ study are the same. Speleology is a cross-disciplinary field that combines the knowledge of chemistry, biology, geology, physics, meteorology, and cartography to develop portraits of caves as complex, evolving systems. History Before modern speleology developed, John Beaumont wrote detailed descriptions of some Mendip caves in the 1680s. The term speleology was coined by Émile Rivière in 1890. Prior to the mid-nineteenth century the scientific value of caves was considered only in its contribution to other branches of science, and cave studies were considered part of the larger disciplines of geography, geology or archaeology. Very little cave-specific study was undertaken prior to the work of Édouard-Alfred Martel (1859–1938), the 'father of modern speleology', who through his extensive and well-publicised cave explorations introduced in France the concept of speleology as a distinct area of study. In 1895 Martel founded the Société de Spéléologie, the first organization devoted to cave science in the world. Other early speleologists include Herbert E. Balch. An international speleological congress was proposed at a meeting in Valence-sur-Rhone, France in 1949 and first held in 1953 in Paris. The International Union of Speleology (UIS) was founded in 1965. The growth of speleology is directly linked with that of the sport of caving, both because of the stimulation of public interest and awareness, and the fact that most speleological field-work has been conducted by sport cavers. Cave geology, hydrogeology and biology Karst is a landscape that has limestone underneath which has been eroded. Caves are usually formed through chemical corrosion via a process of dissolution. Corrosion has several ways of doing this, it can be on carbonate rocks through chemical reactions, in gypsum and rock salt it can happen physically, and in silicate rocks and warm climate the decomposition of the materials can happen as well. Geochemistry Speleothems A speleothem is a geological formation by mineral deposits that accumulate over time in natural caves. Speleothems most commonly form in calcareous caves due to carbonate dissolution reactions. They can take a variety of forms, depending on their depositional history and environment. Their chemical composition, gradual growth, and preservation in caves make them useful paleoclimatic proxies. Biochemistry Caves have an absence of stable temperature, high relative humidity, low rates of evaporation and limited supply of organic material, which help in creating an environment which is highly favorable for the growth of microbes. Microbial assemblages in caves include archaea, bacteria, fungi and other micro-eukaryotes, these highly adapted microbial communities represent the living-backbone of cave ecosystems and play a key role in shaping structures and sustaining trophic networks. Cave cartography The creation of an accurate, detailed map is one of the most common technical activities undertaken within a cave. Cave maps, called surveys, can be used to compare caves to each other by length, depth and volume, may reveal clues on speleogenesis, provide a spatial reference for further scientific study, and assist visitors with route-finding. Cave biology Caves provide a home for many unique biota. Cave ecologies are very diverse, and not sharply distinct from surface habitats. Generally however, the deeper the cave becomes, the more rarefied the ecology. Cave environments fall into three general categories: Endogean the parts of caves that are in communication with surface soils through cracks and rock seams, groundwater seepage, and root protrusion. Parahypogean the threshold regions near cave mouths that extend to the last penetration of sunlight. Hypogean or "true" cave environments. These can be in regular contact with the surface via wind and underground rivers, or the migration of animals, or can be almost entirely isolated. Deep hypogean environments can host autonomous ecologies whose primary source of energy is not sunlight, but chemical energy liberated from limestone and other minerals by chemoautotrophic bacteria. Cave organisms fall into three basic classes: There are also so-called accidental trogloxenes which are surface organisms that enter caves for no survival reason. Some may even be troglophobes (“cave haters”), which cannot survive in caves for any extended period. Examples include deer which fell through a sinkhole, frogs swept into a cave by a flash flood, etc. The two factors that limit cave ecologies are generally energy and nutrients. To some degree moisture is always available in actively forming Karst caves. Cut off from the sunlight and steady deposition of plant detritus, caves are poor habitats in comparison with wet areas on the surface. Most of the energy in cave environments comes from the surplus of the ecosystems outside. One major source of energy and nutrients in caves is dung from trogloxenes, most of which is deposited by bats. Other sources are mentioned above. Cave ecosystems are very fragile. Because of their rarity and position in the ecosystem they are threatened by a large number of human activities. Dam construction, limestone quarrying, water pollution and logging are just some of the disasters that can devastate or destroy underground biological communities. Other areas of cave science Speleologists also work with archaeologists in studying underground ruins, tunnels, sewers and aqueducts, such as the various inlets and outlets of the Cloaca Maxima in Rome.	Physical sciences	Caves	Earth science
152654	https://en.wikipedia.org/wiki/Electrical%20connector	Electrical connector	Components of an electrical circuit are electrically connected if an electric current can run between them through an electrical conductor. An electrical connector is an electromechanical device used to create an electrical connection between parts of an electrical circuit, or between different electrical circuits, thereby joining them into a larger circuit. The connection may be removable (as for portable equipment), require a tool for assembly and removal, or serve as a permanent electrical joint between two points. An adapter can be used to join dissimilar connectors. Most electrical connectors have a genderi.e. the male component, called a plug, connects to the female component, or socket. Thousands of configurations of connectors are manufactured for power, data, and audiovisual applications. Electrical connectors can be divided into four basic categories, differentiated by their function: inline or cable connectors permanently attached to a cable, so it can be plugged into another terminal (either a stationary instrument or another cable) Chassis or panel connectors permanently attached to a piece of equipment so users can connect a cable to a stationary device PCB mount connectors soldered to a printed circuit board, providing a point for cable or wire attachment. (e.g. pin headers, screw terminals, board-to-board connectors) Splice or butt connectors (primarily insulation displacement connectors) that permanently join two lengths of wire or cable In computing, electrical connectors are considered a physical interface and constitute part of the physical layer in the OSI model of networking. Physical construction In addition to the classes mentioned above, connectors are characterised by their pinout, method of connection, materials, size, contact resistance, insulation, mechanical durability, ingress protection, lifetime (number of cycles), and ease of use. It is usually desirable for a connector to be easy to identify visually, rapid to assemble, inexpensive, and require only simple tooling. In some cases an equipment manufacturer might choose a connector specifically because it is not compatible with those from other sources, allowing control of what may be connected. No single connector has all the ideal properties for every application; the proliferation of types is a result of the diverse yet specific requirements of manufacturers. Materials Electrical connectors essentially consist of two classes of materials: conductors and insulators. Properties important to conductor materials are contact resistance, conductivity, mechanical strength, formability, and resilience. Insulators must have a high electrical resistance, withstand high temperatures, and be easy to manufacture for a precise fit Electrodes in connectors are usually made of copper alloys, due to their good conductivity and malleability. Alternatives include brass, phosphor bronze, and beryllium copper. The base electrode metal is often coated with another inert metal such as gold, nickel, or tin. The use of a coating material with good conductivity, mechanical robustness and corrosion resistance helps to reduce the influence of passivating oxide layers and surface adsorbates, which limit metal-to-metal contact patches and contribute to contact resistance. For example, copper alloys have favorable mechanical properties for electrodes, but are hard to solder and prone to corrosion. Thus, copper pins are usually coated with gold to alleviate these pitfalls, especially for analog signals and high-reliability applications. Contact carriers that hold the parts of a connector together are usually made of plastic, due to its insulating properties. Housings or backshells can be made of molded plastic and metal. Connector bodies for high-temperature use, such as thermocouples or associated with large incandescent lamps, may be made of fired ceramic material. Failure modes The majority of connector failures result in intermittent connections or open contacts: Connectors are purely passive componentsthat is, they do not enhance the function of a circuitso connectors should affect the function of a circuit as little as possible. Insecure mounting of connectors (primarily chassis-mounted) can contribute significantly to the risk of failure, especially when subjected to extreme shock or vibration. Other causes of failure are connectors inadequately rated for the applied current and voltage, connectors with inadequate ingress protection, and threaded backshells that are worn or damaged. High temperatures can also cause failure in connectors, resulting in an "avalanche" of failuresambient temperature increases, leading to a decrease in insulation resistance and increase in conductor resistance; this increase generates more heat, and the cycle repeats. Fretting (so-called dynamic corrosion) is a common failure mode in electrical connectors that have not been specifically designed to prevent it, especially in those that are frequently mated and de-mated. Surface corrosion is a risk for many metal parts in connectors, and can cause contacts to form a thin surface layer that increases resistance, thus contributing to heat buildup and intermittent connections. However, remating or reseating a connector can alleviate the issue of surface corrosion, since each cycle scrapes a microscopic layer off the surface of the contact(s), exposing a fresh, unoxidised surface. Circular connectors Many connectors used for industrial and high-reliability applications are circular in cross section, with a cylindrical housing and circular contact interface geometries. This is in contrast to the rectangular design of some connectors, e.g. USB or blade connectors. They are commonly used for easier engagement and disengagement, tight environmental sealing, and rugged mechanical performance. They are widely used in military, aerospace, industrial machinery, and rail, where MIL-DTL-5015 and MIL-DTL-38999 are commonly specified. Fields such as sound engineering and radio communication also use circular connectors, such as XLR and BNC. AC power plugs are also commonly circular, for example, Schuko plugs and IEC 60309. The M12 connector, specified in IEC 61076-2-101, is a circular electrical plug/receptacle pair with 12mm OD mating threads, used in NMEA 2000, DeviceNet, IO-Link, some kinds of Industrial Ethernet, etc. A disadvantage of the circular design is its inefficient use of panel space when used in arrays, when compared to rectangular connectors. Circular connectors commonly use backshells, which provide physical and electromagnetic protection, whilst sometimes also providing a method for locking the connector into a receptacle. In some cases, this backshell provides a hermetic seal, or some degree of ingress protection, through the use of grommets, O-rings, or potting. Hybrid connectors Hybrid connectors allow the intermixing of many connector types, usually by way of a housing with inserts. These housings may also allow intermixing of electrical and non-electrical interfaces, examples of the latter being pneumatic line connectors, and optical fiber connectors. Because hybrid connectors are modular in nature, they tend to simplify assembly, repair, and future modifications. They also allow the creation of composite cable assemblies that can reduce equipment installation time by reducing the number of individual cable and connector assemblies. Mechanical features Pin sequence Some connectors are designed such that certain pins make contact before others when inserted, and break first on disconnection. This is often used in power connectors to protect equipment, e.g. connecting safety ground first. It is also employed for digital signals, as a method to sequence connections properly in hot swapping. Keying Many connectors are keyed with some mechanical component (sometimes called a keyway), which prevents mating in an incorrect orientation. This can be used to prevent mechanical damage to connectors, from being jammed in at the wrong angle or into the wrong connector, or to prevent incompatible or dangerous electrical connections, such as plugging an audio cable into a power outlet. Keying also prevents otherwise symmetrical connectors from being connected in the wrong orientation or polarity. Keying is particularly important for situations where there are many similar connectors, such as in signal electronics. For instance, XLR connectors have a notch to ensure proper orientation, while Mini-DIN plugs have a plastic projection that fits into a corresponding hole in the socket (they also have a notched metal skirt to provide secondary keying). Locking mechanisms Some connector housings are designed with locking mechanisms to prevent inadvertent disconnection or poor environmental sealing. Locking mechanism designs include locking levers of various sorts, jackscrews, screw-in shells, push-pull connector, and toggle or bayonet systems. Some connectors, particularly those with large numbers of contacts, require high forces to connect and disconnect. Locking levers and jackscrews and screw-in shells for such connectors frequently serve both to retain the connector when connected and to provide the force needed for connection and disconnection. Depending on application requirements, housings with locking mechanisms may be tested under various environmental simulations that include physical shock and vibration, water spray, dust, etc. to ensure the integrity of the electrical connection and housing seals. Backshells Backshells are a common accessory for industrial and high-reliability connectors, especially circular connectors. Backshells typically protect the connector and/or cable from environmental or mechanical stress, or shield it from electromagnetic interference. Many types of backshells are available for different purposes, including various sizes, shapes, materials, and levels of protection. Backshells usually lock onto the cable with a clamp or moulded boot, and may be threaded for attachment to a mating receptacle. Backshells for military and aerospace use are regulated by SAE AS85049 within the USA. Hyperboloid contacts To deliver ensured signal stability in extreme environments, traditional pin and socket design may become inadequate. Hyperboloid contacts are designed to withstand more extreme physical demands, such as vibration and shock. They also require around 40% less insertion force as low as per contact,which extends the lifespan, and in some cases offers an alternative to zero insertion force connectors. In a connector with hyperboloid contacts, each female contact has several equally spaced longitudinal wires twisted into a hyperbolic shape. These wires are highly resilient to strain, but still somewhat elastic, hence they essentially function as linear springs. As the male pin is inserted, axial wires in the socket half are deflected, wrapping themselves around the pin to provide a number of contact points. The internal wires that form the hyperboloid structure are usually anchored at each end by bending the tip into a groove or notch in the housing. Whilst hyperboloid contacts may be the only option to make a reliable connection in some circumstances, they have the disadvantage of taking up greater volume in a connector, which can cause problems for high-density connectors. They are also significantly more expensive than traditional pin and socket contacts, which has limited their uptake since their invention in the 1920s by Wilhelm Harold Frederick. In the 1950s, Francois Bonhomme popularised hyperboloid contacts with his "Hypertac" connector, which was later acquired by Smiths Group. During the following decades, the connectors steadily gained popularity, and are still used for medical, industrial, military, aerospace, and rail applications (particularly trains in Europe). Pogo pins Pogo pin or spring loaded connectors are commonly used in consumer and industrial products, where mechanical resilience and ease of use are priorities. The connector consists of a barrel, a spring, and a plunger. They are in applications such as the MagSafe connector where a quick disconnect is desired for safety. Because they rely on spring pressure, not friction, they can be more durable and less damaging than traditional pin and socket design, leading to their use in in-circuit testing. Crown spring connectors Crown spring connectors are commonly used for higher current flows and industrial applications. They have a high number of contact points, which provides a more electrically reliable connection than traditional pin and socket connectors. Methods of connection Whilst technically inaccurate, electrical connectors can be viewed as a type of adapter to convert between two connection methods, which are permanently connected at one end and (usually) detachable at the other end. By definition, each end of this "adapter" has a different connection methode.g. the solder tabs on a male phone connector, and the male phone connector itself. In this example, the solder tabs connected to the cable represent the permanent connection, whilst the male connector portion interfaces with a female socket forming a detachable connection. There are many ways of applying a connector to a cable or device. Some of these methods can be accomplished without specialized tools. Other methods, while requiring a special tool, can assemble connectors much faster and more reliably, and make repairs easier. The number of times a connector can connect and disconnect with its counterpart while meeting all its specifications is termed as mating cycles and is an indirect measure of connector lifespan. The material used for connector contact, plating type and thickness is a major factor that determines the mating cycles. Plug and socket connectors Plug and socket connectors are usually made up of a male plug (typically pin contacts) and a female socket (typically receptacle contacts). Often, but not always, sockets are permanently fixed to a device as in a chassis connector , and plugs are attached to a cable. Plugs generally have one or more pins or prongs that are inserted into openings in the mating socket. The connection between the mating metal parts must be sufficiently tight to make a good electrical connection and complete the circuit. An alternative type of plug and socket connection uses hyperboloid contacts, which makes a more reliable electrical connection. When working with multi-pin connectors, it is helpful to have a pinout diagram to identify the wire or circuit node connected to each pin. Some connector styles may combine pin and socket connection types in a single unit, referred to as a hermaphroditic connector. These connectors includes mating with both male and female aspects, involving complementary paired identical parts each containing both protrusions and indentations. These mating surfaces are mounted into identical fittings that freely mate with any other, without regard for gender (provided that the size and type match). Sometimes both ends of a cable are terminated with the same gender of connector, as in many Ethernet patch cables. In other applications the two ends are terminated differently, either with male and female of the same connector (as in an extension cord), or with incompatible connectors, which is sometimes called an adapter cable. Plugs and sockets are widely used in various connector systems including blade connectors, breadboards, XLR connectors, car power outlets, banana connectors, and phone connectors. Jacks and plugs A jack is a connector that installs on the surface of a bulkhead or enclosure, and mates with its reciprocal, the plug. According to the American Society of Mechanical Engineers, the stationary (more fixed) connector of a pair is classified as a jack (denoted J), usually attached to a piece of equipment as in a chassis-mount or panel-mount connector. The movable (less fixed) connector is classified as a plug (denoted P), designed to attach to a wire, cable or removable electrical assembly. This convention is currently defined in ASME Y14.44-2008, which supersedes IEEE 200-1975, which in turn derives from the long-withdrawn MIL-STD-16 (from the 1950s), highlighting the heritage of this connector naming convention. IEEE 315-1975 works alongside ASME Y14.44-2008 to define jacks and plugs. The term jack occurs in several related terms: The registered jack or modular jack in RJ11, RJ45 and other similar connectors used for telecommunications and computer networking The telephone jack of manual telephone switchboards, which is the socket fitting the original telephone plug The phone jack common to many electronic applications in various configurations, sometimes referred to as a headphone jack The RCA jack, also known as a phono jack, common to consumer audiovisual electronics The EIAJ jack for consumer appliances requiring a power supply of less than 18.0 volts Crimp-on connectors Crimped connectors are a type of solderless connection, using mechanical friction and uniform deformation to secure a connector to a pre-stripped wire (usually stranded). Crimping is used in splice connectors, crimped multipin plugs and sockets, and crimped coaxial connectors. Crimping usually requires a specialised crimping tool, but the connectors are quick and easy to install and are a common alternative to solder connections or insulation displacement connectors. Effective crimp connections deform the metal of the connector past its yield point so that the compressed wire causes tension in the surrounding connector, and these forces counter each other to create a high degree of static friction. Due to the elastic element in crimped connections, they are highly resistant to vibration and thermal shock. Crimped contacts are permanent (i.e. the connectors and wire ends cannot be reused). Crimped plug-and-socket connectors can be classified as rear release or front release. This relates to the side of the connector where the pins are anchored: Front release contacts are released from the front (contact side) of the connector, and removed from the rear. The removal tool engages with the front portion of the contact and pushes it through to the back of the connector. Rear release contacts are released and removed from the rear (wire side) of the connector. The removal tool releases the contacts from the rear and pulls the contact out of the retainer. Soldered connectors Many plug and socket connectors are attached to a wire or cable by soldering conductors to electrodes on the back of the connector. Soldered joints in connectors are robust and reliable if executed correctly, but are usually slower to make than crimped connections. When wires are to be soldered to the back of a connector, a backshell is often used to protect the connection and add strain relief. Metal solder buckets or solder cups are provided, which consist of a cylindrical cavity that an installer fills with solder before inserting the wire. When creating soldered connections, it is possible to melt the dielectric between pins or wires. This can cause problems because the thermal conductivity of metals causes heat to quickly distribute through the cable and connector, and when this heat melts plastic dielectric, it can cause short circuits or "flared" (conical) insulation. Solder joints are also more prone to mechanical failure than crimped joints when subjected to vibration and compression. Insulation-displacement connectors Since stripping insulation from wires is time-consuming, many connectors intended for rapid assembly use insulation-displacement connectors which cut the insulation as the wire is inserted. These generally take the form of a fork-shaped opening in the terminal, into which the insulated wire is pressed, which cut through the insulation to contact the conductor. To make these connections reliably on a production line, special tools accurately control the forces applied during assembly. On small scales, these tools tend to cost more than tools for crimped connections. Insulation displacement connectors are usually used with small conductors for signal purposes and at low voltage. Power conductors carrying more than a few amperes are more reliably terminated with other means, though "hot tap" press-on connectors find some use in automotive applications for additions to existing wiring. A common example is the multi-conductor flat ribbon cable used in computer disk drives; to terminate each of the many (approximately 40) wires individually would be slow and error-prone, but an insulation displacement connector can terminate all the wires in a single action. Another very common use is so-called punch-down blocks used for terminating unshielded twisted pair wiring. Binding posts Binding posts are a single-wire connection method, where stripped wire is screwed or clamped to a metal electrode. Such connectors are frequently used in electronic test equipment and audio. Many binding posts also accept a banana plug. Screw terminals Screw connections are frequently used for semi-permanent wiring and connections inside devices, due to their simple but reliable construction. The basic principle of all screw terminals involves the tip of a bolt clamping onto a stripped conductor. They can be used to join multiple conductors, to connect wires to a printed circuit board, or to terminate a cable into a plug or socket. The clamping screw may act in the longitudinal axis (parallel to the wire) or the transverse axis (perpendicular to the wire), or both. Some disadvantages are that connecting wires is more difficult than simply plugging in a cable, and screw terminals are generally not very well protected from contact with persons or foreign conducting materials. Terminal blocks (also called terminal boards or strips) provide a convenient means of connecting individual electrical wires without a splice or physically joining the ends. Since terminal blocks are readily available for a wide range of wire sizes and terminal quantity, they are one of the most flexible types of electrical connector available. One type of terminal block accepts wires that are prepared only by stripping a short length of insulation from the end. Another type, often called barrier strips, accepts wires that have ring or spade terminal lugs crimped onto the wires. Printed circuit board (PCB) mounted screw terminals let individual wires connect to a PCB through leads soldered to the board. Ring and spade connectors The connectors in the top row of the image are known as ring terminals and spade terminals (sometimes called fork or split ring terminals). Electrical contact is made by the flat surface of the ring or spade, while mechanically they are attached by passing a screw or bolt through them. The spade terminal form factor facilitates connections since the screw or bolt can be left partially screwed in as the spade terminal is removed or attached. Their sizes can be determined by the gauge of the conducting wire, and the interior and exterior diameters. In the case of insulated crimp connectors, the crimped area lies under an insulating sleeve through which the pressing force acts. During crimping, the extended end of this insulating sleeve is simultaneously pressed around the insulated area of the cable, creating strain relief. The insulating sleeve of insulated connectors has a color that indicates the wire's cross-section area. Colors are standardized according to DIN 46245: Red for cross-section areas from 0.5 to 1 mm² Blue for cross-section areas from 1.5 to 2.5 mm² Yellow for cross-section areas over 4 to 6 mm² Blade connectors A blade connector is a type of single wire, plug-and-socket connection device using a flat conductive blade (plug) that is inserted into a receptacle. Wires are typically attached to male or female blade connector terminals by either crimping or soldering. Insulated and uninsulated varieties are available. In some cases the blade is an integral manufactured part of a component (such as a switch or a speaker unit), and the reciprocal connector terminal is pushed onto the device's connector terminal. Other connection methods Alligator and Crocodile clips – conductive clamps used for temporary connections, e.g. jumper cables Board to board connectors – e.g. card-edge connectors or FPGA mezzanine connectors Twist-on wire connectors (e.g. wire nuts) – used in low-voltage power circuits for wires up to about 10 AWG Wire wrapping – used in older circuit boards	Technology	Components	null
152664	https://en.wikipedia.org/wiki/Loop%20quantum%20gravity	Loop quantum gravity	Loop quantum gravity (LQG) is a theory of quantum gravity that incorporates matter of the Standard Model into the framework established for the intrinsic quantum gravity case. It is an attempt to develop a quantum theory of gravity based directly on Albert Einstein's geometric formulation rather than the treatment of gravity as a mysterious mechanism (force). As a theory, LQG postulates that the structure of space and time is composed of finite loops woven into an extremely fine fabric or network. These networks of loops are called spin networks. The evolution of a spin network, or spin foam, has a scale on the order of a Planck length, approximately 10−35 meters, and smaller scales are meaningless. Consequently, not just matter, but space itself, prefers an atomic structure. The areas of research, which involve about 30 research groups worldwide, share the basic physical assumptions and the mathematical description of quantum space. Research has evolved in two directions: the more traditional canonical loop quantum gravity, and the newer covariant loop quantum gravity, called spin foam theory. The most well-developed theory that has been advanced as a direct result of loop quantum gravity is called loop quantum cosmology (LQC). LQC advances the study of the early universe, incorporating the concept of the Big Bang into the broader theory of the Big Bounce, which envisions the Big Bang as the beginning of a period of expansion, that follows a period of contraction, which has been described as the Big Crunch. History In 1986, Abhay Ashtekar reformulated Einstein's general relativity in a language closer to that of the rest of fundamental physics, specifically Yang–Mills theory. Shortly after, Ted Jacobson and Lee Smolin realized that the formal equation of quantum gravity, called the Wheeler–DeWitt equation, admitted solutions labelled by loops when rewritten in the new Ashtekar variables. Carlo Rovelli and Smolin defined a nonperturbative and background-independent quantum theory of gravity in terms of these loop solutions. Jorge Pullin and Jerzy Lewandowski understood that the intersections of the loops are essential for the consistency of the theory, and the theory should be formulated in terms of intersecting loops, or graphs. In 1994, Rovelli and Smolin showed that the quantum operators of the theory associated to area and volume have a discrete spectrum. That is, geometry is quantized. This result defines an explicit basis of states of quantum geometry, which turned out to be labelled by Roger Penrose's spin networks, which are graphs labelled by spins. The canonical version of the dynamics was established by Thomas Thiemann, who defined an anomaly-free Hamiltonian operator and showed the existence of a mathematically consistent background-independent theory. The covariant, or "spin foam", version of the dynamics was developed jointly over several decades by research groups in France, Canada, UK, Poland, and Germany. It was completed in 2008, leading to the definition of a family of transition amplitudes, which in the classical limit can be shown to be related to a family of truncations of general relativity. The finiteness of these amplitudes was proven in 2011. It requires the existence of a positive cosmological constant, which is consistent with observed acceleration in the expansion of the Universe. Background independence LQG is formally background independent, meaning the equations of LQG are not embedded in, or dependent on, space and time (except for its invariant topology). Instead, they are expected to give rise to space and time at distances which are 10 times the Planck length. The issue of background independence in LQG still has some unresolved subtleties. For example, some derivations require a fixed choice of the topology, while any consistent quantum theory of gravity should include topology change as a dynamical process. Spacetime as a "container" over which physics takes place has no objective physical meaning and instead the gravitational interaction is represented as just one of the fields forming the world. This is known as the relationalist interpretation of spacetime. In LQG this aspect of general relativity is taken seriously and this symmetry is preserved by requiring that the physical states remain invariant under the generators of diffeomorphisms. The interpretation of this condition is well understood for purely spatial diffeomorphisms. However, the understanding of diffeomorphisms involving time (the Hamiltonian constraint) is more subtle because it is related to dynamics and the so-called "problem of time" in general relativity. A generally accepted calculational framework to account for this constraint has yet to be found. A plausible candidate for the quantum Hamiltonian constraint is the operator introduced by Thiemann. Constraints and their Poisson bracket algebra Dirac observables The constraints define a constraint surface in the original phase space. The gauge motions of the constraints apply to all phase space but have the feature that they leave the constraint surface where it is, and thus the orbit of a point in the hypersurface under gauge transformations will be an orbit entirely within it. Dirac observables are defined as phase space functions, , that Poisson commute with all the constraints when the constraint equations are imposed, that is, they are quantities defined on the constraint surface that are invariant under the gauge transformations of the theory. Then, solving only the constraint and determining the Dirac observables with respect to it leads us back to the Arnowitt–Deser–Misner (ADM) phase space with constraints . The dynamics of general relativity is generated by the constraints, it can be shown that six Einstein equations describing time evolution (really a gauge transformation) can be obtained by calculating the Poisson brackets of the three-metric and its conjugate momentum with a linear combination of the spatial diffeomorphism and Hamiltonian constraint. The vanishing of the constraints, giving the physical phase space, are the four other Einstein equations. Quantization of the constraints – the equations of quantum general relativity Pre-history and Ashtekar new variables Many of the technical problems in canonical quantum gravity revolve around the constraints. Canonical general relativity was originally formulated in terms of metric variables, but there seemed to be insurmountable mathematical difficulties in promoting the constraints to quantum operators because of their highly non-linear dependence on the canonical variables. The equations were much simplified with the introduction of Ashtekar's new variables. Ashtekar variables describe canonical general relativity in terms of a new pair of canonical variables closer to those of gauge theories. The first step consists of using densitized triads (a triad is simply three orthogonal vector fields labeled by and the densitized triad is defined by ) to encode information about the spatial metric, (where is the flat space metric, and the above equation expresses that , when written in terms of the basis , is locally flat). (Formulating general relativity with triads instead of metrics was not new.) The densitized triads are not unique, and in fact one can perform a local in space rotation with respect to the internal indices . The canonically conjugate variable is related to the extrinsic curvature by . But problems similar to using the metric formulation arise when one tries to quantize the theory. Ashtekar's new insight was to introduce a new configuration variable, that behaves as a complex connection where is related to the so-called spin connection via . Here is called the chiral spin connection. It defines a covariant derivative . It turns out that is the conjugate momentum of , and together these form Ashtekar's new variables. The expressions for the constraints in Ashtekar variables; Gauss's theorem, the spatial diffeomorphism constraint and the (densitized) Hamiltonian constraint then read: respectively, where is the field strength tensor of the connection and where is referred to as the vector constraint. The above-mentioned local in space rotational invariance is the original of the gauge invariance here expressed by Gauss's theorem. Note that these constraints are polynomial in the fundamental variables, unlike the constraints in the metric formulation. This dramatic simplification seemed to open up the way to quantizing the constraints. (See the article Self-dual Palatini action for a derivation of Ashtekar's formalism). With Ashtekar's new variables, given the configuration variable , it is natural to consider wavefunctions . This is the connection representation. It is analogous to ordinary quantum mechanics with configuration variable and wavefunctions . The configuration variable gets promoted to a quantum operator via: (analogous to ) and the triads are (functional) derivatives, (analogous to ). In passing over to the quantum theory the constraints become operators on a kinematic Hilbert space (the unconstrained Yang–Mills Hilbert space). Note that different ordering of the 's and 's when replacing the 's with derivatives give rise to different operators – the choice made is called the factor ordering and should be chosen via physical reasoning. Formally they read There are still problems in properly defining all these equations and solving them. For example, the Hamiltonian constraint Ashtekar worked with was the densitized version instead of the original Hamiltonian, that is, he worked with . There were serious difficulties in promoting this quantity to a quantum operator. Moreover, although Ashtekar variables had the virtue of simplifying the Hamiltonian, they are complex. When one quantizes the theory, it is difficult to ensure that one recovers real general relativity as opposed to complex general relativity. Quantum constraints as the equations of quantum general relativity The classical result of the Poisson bracket of the smeared Gauss' law with the connections is The quantum Gauss' law reads If one smears the quantum Gauss' law and study its action on the quantum state one finds that the action of the constraint on the quantum state is equivalent to shifting the argument of by an infinitesimal (in the sense of the parameter small) gauge transformation, and the last identity comes from the fact that the constraint annihilates the state. So the constraint, as a quantum operator, is imposing the same symmetry that its vanishing imposed classically: it is telling us that the functions have to be gauge invariant functions of the connection. The same idea is true for the other constraints. Therefore, the two step process in the classical theory of solving the constraints (equivalent to solving the admissibility conditions for the initial data) and looking for the gauge orbits (solving the 'evolution' equations) is replaced by a one step process in the quantum theory, namely looking for solutions of the quantum equations . This is because it solves the constraint at the quantum level and it simultaneously looks for states that are gauge invariant because is the quantum generator of gauge transformations (gauge invariant functions are constant along the gauge orbits and thus characterize them). Recall that, at the classical level, solving the admissibility conditions and evolution equations was equivalent to solving all of Einstein's field equations, this underlines the central role of the quantum constraint equations in canonical quantum gravity. Introduction of the loop representation It was in particular the inability to have good control over the space of solutions to Gauss's law and spatial diffeomorphism constraints that led Rovelli and Smolin to consider the loop representation in gauge theories and quantum gravity. LQG includes the concept of a holonomy. A holonomy is a measure of how much the initial and final values of a spinor or vector differ after parallel transport around a closed loop; it is denoted . Knowledge of the holonomies is equivalent to knowledge of the connection, up to gauge equivalence. Holonomies can also be associated with an edge; under a Gauss Law these transform as For a closed loop and assuming , yields or The trace of an holonomy around a closed loop is written and is called a Wilson loop. Thus Wilson loops are gauge invariant. The explicit form of the Holonomy is where is the curve along which the holonomy is evaluated, and is a parameter along the curve, denotes path ordering meaning factors for smaller values of appear to the left, and are matrices that satisfy the algebra The Pauli matrices satisfy the above relation. It turns out that there are infinitely many more examples of sets of matrices that satisfy these relations, where each set comprises matrices with , and where none of these can be thought to 'decompose' into two or more examples of lower dimension. They are called different irreducible representations of the algebra. The most fundamental representation being the Pauli matrices. The holonomy is labelled by a half integer according to the irreducible representation used. The use of Wilson loops explicitly solves the Gauss gauge constraint. Loop representation is required to handle the spatial diffeomorphism constraint. With Wilson loops as a basis, any Gauss gauge invariant function expands as, This is called the loop transform and is analogous to the momentum representation in quantum mechanics (see Position and momentum space). The QM representation has a basis of states labelled by a number and expands as and works with the coefficients of the expansion The inverse loop transform is defined by This defines the loop representation. Given an operator in the connection representation, one should define the corresponding operator on in the loop representation via, where is defined by the usual inverse loop transform, A transformation formula giving the action of the operator on in terms of the action of the operator on is then obtained by equating the R.H.S. of with the R.H.S. of with substituted into , namely or where means the operator but with the reverse factor ordering (remember from simple quantum mechanics where the product of operators is reversed under conjugation). The action of this operator on the Wilson loop is evaluated as a calculation in the connection representation and the result is rearranged purely as a manipulation in terms of loops (with regard to the action on the Wilson loop, the chosen transformed operator is the one with the opposite factor ordering compared to the one used for its action on wavefunctions ). This gives the physical meaning of the operator . For example, if corresponded to a spatial diffeomorphism, then this can be thought of as keeping the connection field of where it is while performing a spatial diffeomorphism on instead. Therefore, the meaning of is a spatial diffeomorphism on , the argument of . In the loop representation, the spatial diffeomorphism constraint is solved by considering functions of loops that are invariant under spatial diffeomorphisms of the loop . That is, knot invariants are used. This opens up an unexpected connection between knot theory and quantum gravity. Any collection of non-intersecting Wilson loops satisfy Ashtekar's quantum Hamiltonian constraint. Using a particular ordering of terms and replacing by a derivative, the action of the quantum Hamiltonian constraint on a Wilson loop is When a derivative is taken it brings down the tangent vector, , of the loop, . So, However, as is anti-symmetric in the indices and this vanishes (this assumes that is not discontinuous anywhere and so the tangent vector is unique). With regard to loop representation, the wavefunctions vanish when the loop has discontinuities and are knot invariants. Such functions solve the Gauss law, the spatial diffeomorphism constraint and (formally) the Hamiltonian constraint. This yields an infinite set of exact (if only formal) solutions to all the equations of quantum general relativity! This generated a lot of interest in the approach and eventually led to LQG. Geometric operators, the need for intersecting Wilson loops and spin network states The easiest geometric quantity is the area. Let us choose coordinates so that the surface is characterized by . The area of small parallelogram of the surface is the product of length of each side times where is the angle between the sides. Say one edge is given by the vector and the other by then, In the space spanned by and there is an infinitesimal parallelogram described by and . Using (where the indices and run from 1 to 2), yields the area of the surface given by where and is the determinant of the metric induced on . The latter can be rewritten where the indices go from 1 to 2. This can be further rewritten as The standard formula for an inverse matrix is There is a similarity between this and the expression for . But in Ashtekar variables, . Therefore, According to the rules of canonical quantization the triads should be promoted to quantum operators, The area can be promoted to a well defined quantum operator despite the fact that it contains a product of two functional derivatives and a square-root. Putting (-th representation), This quantity is important in the final formula for the area spectrum. The result is where the sum is over all edges of the Wilson loop that pierce the surface . The formula for the volume of a region is given by The quantization of the volume proceeds the same way as with the area. Each time the derivative is taken, it brings down the tangent vector , and when the volume operator acts on non-intersecting Wilson loops the result vanishes. Quantum states with non-zero volume must therefore involve intersections. Given that the anti-symmetric summation is taken over in the formula for the volume, it needs intersections with at least three non-coplanar lines. At least four-valent vertices are needed for the volume operator to be non-vanishing. Assuming the real representation where the gauge group is , Wilson loops are an over complete basis as there are identities relating different Wilson loops. These occur because Wilson loops are based on matrices (the holonomy) and these matrices satisfy identities. Given any two matrices and , This implies that given two loops and that intersect, where by we mean the loop traversed in the opposite direction and means the loop obtained by going around the loop and then along . See figure below. Given that the matrices are unitary one has that . Also given the cyclic property of the matrix traces (i.e. ) one has that . These identities can be combined with each other into further identities of increasing complexity adding more loops. These identities are the so-called Mandelstam identities. Spin networks certain are linear combinations of intersecting Wilson loops designed to address the over-completeness introduced by the Mandelstam identities (for trivalent intersections they eliminate the over-completeness entirely) and actually constitute a basis for all gauge invariant functions. As mentioned above the holonomy tells one how to propagate test spin half particles. A spin network state assigns an amplitude to a set of spin half particles tracing out a path in space, merging and splitting. These are described by spin networks : the edges are labelled by spins together with 'intertwiners' at the vertices which are prescription for how to sum over different ways the spins are rerouted. The sum over rerouting are chosen as such to make the form of the intertwiner invariant under Gauss gauge transformations. Hamiltonian constraint of LQG In the long history of canonical quantum gravity formulating the Hamiltonian constraint as a quantum operator (Wheeler–DeWitt equation) in a mathematically rigorous manner has been a formidable problem. It was in the loop representation that a mathematically well defined Hamiltonian constraint was finally formulated in 1996. We leave more details of its construction to the article Hamiltonian constraint of LQG. This together with the quantum versions of the Gauss law and spatial diffeomorphism constrains written in the loop representation are the central equations of LQG (modern canonical quantum General relativity). Finding the states that are annihilated by these constraints (the physical states), and finding the corresponding physical inner product, and observables is the main goal of the technical side of LQG. An important aspect of the Hamiltonian operator is that it only acts at vertices (a consequence of this is that Thiemann's Hamiltonian operator, like Ashtekar's operator, annihilates non-intersecting loops except now it is not just formal and has rigorous mathematical meaning). More precisely, its action is non-zero on at least vertices of valence three and greater and results in a linear combination of new spin networks where the original graph has been modified by the addition of lines at each vertex together and a change in the labels of the adjacent links of the vertex. Chiral fermions and the fermion doubling problem A significant challenge in theoretical physics lies in unifying LQG, a theory of quantum spacetime, with the Standard Model of particle physics, which describes fundamental forces and particles. A major obstacle in this endeavor is the fermion doubling problem, which arises when incorporating chiral fermions into the LQG framework. Chiral fermions, such as electrons and quarks, are fundamental particles characterized by their "handedness" or chirality. This property dictates that a particle and its mirror image behave differently under weak interactions. This asymmetry is fundamental to the Standard Model's success in explaining numerous physical phenomena. However, attempts to integrate chiral fermions into LQG often result in the appearance of spurious, mirror-image particles. Instead of a single left-handed fermion, for instance, the theory predicts the existence of both a left-handed and a right-handed version. This "doubling" contradicts the observed chirality of the Standard Model and disrupts its predictive power. The fermion doubling problem poses a significant hurdle in constructing a consistent theory of quantum gravity. The Standard Model's accuracy in describing the universe at the smallest scales relies heavily on the unique properties of chiral fermions. Without a solution to this problem, incorporating matter and its interactions into a unified framework of quantum gravity remains a significant challenge. Therefore, resolving the fermion doubling problem is crucial for advancing our understanding of the universe at its most fundamental level and developing a complete theory that unites gravity with the quantum world. Spin foams In loop quantum gravity (LQG), a spin network represents a "quantum state" of the gravitational field on a 3-dimensional hypersurface. The set of all possible spin networks (or, more accurately, "s-knots" – that is, equivalence classes of spin networks under diffeomorphisms) is countable; it constitutes a basis of LQG Hilbert space. In physics, a spin foam is a topological structure made out of two-dimensional faces that represents one of the configurations that must be summed to obtain a Feynman's path integral (functional integration) description of quantum gravity. It is closely related to loop quantum gravity. Spin foam derived from the Hamiltonian constraint operator On this section see and references therein. The Hamiltonian constraint generates 'time' evolution. Solving the Hamiltonian constraint should tell us how quantum states evolve in 'time' from an initial spin network state to a final spin network state. One approach to solving the Hamiltonian constraint starts with what is called the Dirac delta function. The summation of which over different sequences of actions can be visualized as a summation over different histories of 'interaction vertices' in the 'time' evolution sending the initial spin network to the final spin network. Each time a Hamiltonian operator acts it does so by adding a new edge at the vertex. This then naturally gives rise to the two-complex (a combinatorial set of faces that join along edges, which in turn join on vertices) underlying the spin foam description; we evolve forward an initial spin network sweeping out a surface, the action of the Hamiltonian constraint operator is to produce a new planar surface starting at the vertex. We are able to use the action of the Hamiltonian constraint on the vertex of a spin network state to associate an amplitude to each "interaction" (in analogy to Feynman diagrams). See figure below. This opens a way of trying to directly link canonical LQG to a path integral description. Just as a spin networks describe quantum space, each configuration contributing to these path integrals, or sums over history, describe 'quantum spacetime'. Because of their resemblance to soap foams and the way they are labeled John Baez gave these 'quantum spacetimes' the name 'spin foams'. There are however severe difficulties with this particular approach, for example the Hamiltonian operator is not self-adjoint, in fact it is not even a normal operator (i.e. the operator does not commute with its adjoint) and so the spectral theorem cannot be used to define the exponential in general. The most serious problem is that the 's are not mutually commuting, it can then be shown the formal quantity cannot even define a (generalized) projector. The master constraint (see below) does not suffer from these problems and as such offers a way of connecting the canonical theory to the path integral formulation. Spin foams from BF theory It turns out there are alternative routes to formulating the path integral, however their connection to the Hamiltonian formalism is less clear. One way is to start with the BF theory. This is a simpler theory than general relativity, it has no local degrees of freedom and as such depends only on topological aspects of the fields. BF theory is what is known as a topological field theory. Surprisingly, it turns out that general relativity can be obtained from BF theory by imposing a constraint, BF theory involves a field and if one chooses the field to be the (anti-symmetric) product of two tetrads (tetrads are like triads but in four spacetime dimensions), one recovers general relativity. The condition that the field be given by the product of two tetrads is called the simplicity constraint. The spin foam dynamics of the topological field theory is well understood. Given the spin foam 'interaction' amplitudes for this simple theory, one then tries to implement the simplicity conditions to obtain a path integral for general relativity. The non-trivial task of constructing a spin foam model is then reduced to the question of how this simplicity constraint should be imposed in the quantum theory. The first attempt at this was the famous Barrett–Crane model. However this model was shown to be problematic, for example there did not seem to be enough degrees of freedom to ensure the correct classical limit. It has been argued that the simplicity constraint was imposed too strongly at the quantum level and should only be imposed in the sense of expectation values just as with the Lorenz gauge condition in the Gupta–Bleuler formalism of quantum electrodynamics. New models have now been put forward, sometimes motivated by imposing the simplicity conditions in a weaker sense. Another difficulty here is that spin foams are defined on a discretization of spacetime. While this presents no problems for a topological field theory as it has no local degrees of freedom, it presents problems for GR. This is known as the problem triangularization dependence. Modern formulation of spin foams Just as imposing the classical simplicity constraint recovers general relativity from BF theory, it is expected that an appropriate quantum simplicity constraint will recover quantum gravity from quantum BF theory. Progress has been made with regard to this issue by Engle, Pereira, and Rovelli, Freidel and Krasnov and Livine and Speziale in defining spin foam interaction amplitudes with better behaviour. An attempt to make contact between EPRL-FK spin foam and the canonical formulation of LQG has been made. Spin foam derived from the master constraint operator See below. The semiclassical limit and loop quantum gravity The Classical limit is the ability of a physical theory to approximate classical mechanics. It is used with physical theories that predict non-classical behavior. Any candidate theory of quantum gravity must be able to reproduce Einstein's theory of general relativity as a classical limit of a quantum theory. This is not guaranteed because of a feature of quantum field theories which is that they have different sectors, these are analogous to the different phases that come about in the thermodynamical limit of statistical systems. Just as different phases are physically different, so are different sectors of a quantum field theory. It may turn out that LQG belongs to an unphysical sector – one in which one does not recover general relativity in the semiclassical limit or there might not be any physical sector. Moreover, the physical Hilbert space must contain enough semiclassical states to guarantee that the quantum theory obtained can return to the classical theory when avoiding quantum anomalies; otherwise there will be restrictions on the physical Hilbert space that have no counterpart in the classical theory, implying that the quantum theory has fewer degrees of freedom than the classical theory. Theorems establishing the uniqueness of the loop representation as defined by Ashtekar et al. (i.e. a certain concrete realization of a Hilbert space and associated operators reproducing the correct loop algebra) have been given by two groups (Lewandowski, Okołów, Sahlmann and Thiemann; and Christian Fleischhack). Before this result was established it was not known whether there could be other examples of Hilbert spaces with operators invoking the same loop algebra – other realizations not equivalent to the one that had been used. These uniqueness theorems imply no others exist, so if LQG does not have the correct semiclassical limit then the theorems would mean the end of the loop representation of quantum gravity. Difficulties and progress checking the semiclassical limit There are a number of difficulties in trying to establish LQG gives Einstein's theory of general relativity in the semiclassical limit: There is no operator corresponding to infinitesimal spatial diffeomorphisms (it is not surprising that the theory has no generator of infinitesimal spatial 'translations' as it predicts spatial geometry has a discrete nature, compare to the situation in condensed matter). Instead it must be approximated by finite spatial diffeomorphisms and so the Poisson bracket structure of the classical theory is not exactly reproduced. This problem can be circumvented with the introduction of the so-called master constraint (see below). There is the problem of reconciling the discrete combinatorial nature of the quantum states with the continuous nature of the fields of the classical theory. There are serious difficulties arising from the structure of the Poisson brackets involving the spatial diffeomorphism and Hamiltonian constraints. In particular, the algebra of (smeared) Hamiltonian constraints does not close: It is proportional to a sum over infinitesimal spatial diffeomorphisms (which, as noted above, does not exist in the quantum theory) where the coefficients of proportionality are not constants but have non-trivial phase space dependence – as such it does not form a Lie algebra. However, the situation is improved by the introduction of the master constraint. The semiclassical machinery developed so far is only appropriate to non-graph-changing operators, however, Thiemann's Hamiltonian constraint is a graph-changing operator – the new graph it generates has degrees of freedom upon which the coherent state does not depend and so their quantum fluctuations are not suppressed. There is also the restriction, so far, that these coherent states are only defined at the Kinematic level, and now one has to lift them to the level of and . It can be shown that Thiemann's Hamiltonian constraint is required to be graph-changing in order to resolve problem 3 in some sense. The master constraint algebra however is trivial and so the requirement that it be graph-changing can be lifted and indeed non-graph-changing master constraint operators have been defined. As far as is currently known, this problem is still out of reach. Formulating observables for classical general relativity is a formidable problem because of its non-linear nature and spacetime diffeomorphism invariance. A systematic approximation scheme to calculate observables has been recently developed. Difficulties in trying to examine the semiclassical limit of the theory should not be confused with it having the wrong semiclassical limit. Concerning issue number 2 above, consider so-called weave states. Ordinary measurements of geometric quantities are macroscopic, and Planckian discreteness is smoothed out. The fabric of a T-shirt is analogous: at a distance it is a smooth curved two-dimensional surface, but on closer inspection we see that it is actually composed of thousands of one-dimensional linked threads. The image of space given in LQG is similar. Consider a large spin network formed by a large number of nodes and links, each of Planck scale. Probed at a macroscopic scale, it appears as a three-dimensional continuous metric geometry. To make contact with low energy physics it is mandatory to develop approximation schemes both for the physical inner product and for Dirac observables; the spin foam models that have been intensively studied can be viewed as avenues toward approximation schemes for said physical inner product. Markopoulou, et al. adopted the idea of noiseless subsystems in an attempt to solve the problem of the low energy limit in background independent quantum gravity theories. The idea has led to the possibility of matter of the standard model being identified with emergent degrees of freedom from some versions of LQG (see section below: LQG and related research programs). As Wightman emphasized in the 1950s, in Minkowski QFTs the point functions completely determine the theory. In particular, one can calculate the scattering amplitudes from these quantities. As explained below in the section on the Background independent scattering amplitudes, in the background-independent context, the point functions refer to a state and in gravity that state can naturally encode information about a specific geometry which can then appear in the expressions of these quantities. To leading order, LQG calculations have been shown to agree in an appropriate sense with the point functions calculated in the effective low energy quantum general relativity. Improved dynamics and the master constraint The master constraint Thiemann's Master Constraint Programme for Loop Quantum Gravity (LQG) was proposed as a classically equivalent way to impose the infinite number of Hamiltonian constraint equations in terms of a single master constraint , which involves the square of the constraints in question. An initial objection to the use of the master constraint was that on first sight it did not seem to encode information about the observables; because the Master constraint is quadratic in the constraint, when one computes its Poisson bracket with any quantity, the result is proportional to the constraint, therefore it vanishes when the constraints are imposed and as such does not select out particular phase space functions. However, it was realized that the condition is where is at least a twice differentiable function on phase space is equivalent to being a weak Dirac observable with respect to the constraints in question. So the master constraint does capture information about the observables. Because of its significance this is known as the master equation. That the master constraint Poisson algebra is an honest Lie algebra opens the possibility of using a method, known as group averaging, in order to construct solutions of the infinite number of Hamiltonian constraints, a physical inner product thereon and Dirac observables via what is known as refined algebraic quantization, or RAQ. The quantum master constraint Define the quantum master constraint (regularisation issues aside) as Obviously, for all implies . Conversely, if then implies . First compute the matrix elements of the would-be operator , that is, the quadratic form . is a graph changing, diffeomorphism invariant quadratic form that cannot exist on the kinematic Hilbert space , and must be defined on . Since the master constraint operator is densely defined on , then is a positive and symmetric operator in . Therefore, the quadratic form associated with is closable. The closure of is the quadratic form of a unique self-adjoint operator , called the Friedrichs extension of . We relabel as for simplicity. Note that the presence of an inner product, viz Eq 4, means there are no superfluous solutions i.e. there are no such that but for which . It is also possible to construct a quadratic form for what is called the extended master constraint (discussed below) on which also involves the weighted integral of the square of the spatial diffeomorphism constraint (this is possible because is not graph changing). The spectrum of the master constraint may not contain zero due to normal or factor ordering effects which are finite but similar in nature to the infinite vacuum energies of background-dependent quantum field theories. In this case it turns out to be physically correct to replace with provided that the "normal ordering constant" vanishes in the classical limit, that is, so that is a valid quantisation of . Testing the master constraint The constraints in their primitive form are rather singular, this was the reason for integrating them over test functions to obtain smeared constraints. However, it would appear that the equation for the master constraint, given above, is even more singular involving the product of two primitive constraints (although integrated over space). Squaring the constraint is dangerous as it could lead to worsened ultraviolet behaviour of the corresponding operator and hence the master constraint programme must be approached with care. In doing so the master constraint programme has been satisfactorily tested in a number of model systems with non-trivial constraint algebras, free and interacting field theories. The master constraint for LQG was established as a genuine positive self-adjoint operator and the physical Hilbert space of LQG was shown to be non-empty, a consistency test LQG must pass to be a viable theory of quantum general relativity. Applications of the master constraint The master constraint has been employed in attempts to approximate the physical inner product and define more rigorous path integrals. The Consistent Discretizations approach to LQG, is an application of the master constraint program to construct the physical Hilbert space of the canonical theory. Spin foam from the master constraint The master constraint is easily generalized to incorporate the other constraints. It is then referred to as the extended master constraint, denoted . We can define the extended master constraint which imposes both the Hamiltonian constraint and spatial diffeomorphism constraint as a single operator, . Setting this single constraint to zero is equivalent to and for all in . This constraint implements the spatial diffeomorphism and Hamiltonian constraint at the same time on the Kinematic Hilbert space. The physical inner product is then defined as (as ). A spin foam representation of this expression is obtained by splitting the -parameter in discrete steps and writing The spin foam description then follows from the application of on a spin network resulting in a linear combination of new spin networks whose graph and labels have been modified. Obviously an approximation is made by truncating the value of to some finite integer. An advantage of the extended master constraint is that we are working at the kinematic level and so far it is only here we have access semiclassical coherent states. Moreover, one can find none graph changing versions of this master constraint operator, which are the only type of operators appropriate for these coherent states. Algebraic quantum gravity (AQG) The master constraint programme has evolved into a fully combinatorial treatment of gravity known as algebraic quantum gravity (AQG). The non-graph changing master constraint operator is adapted in the framework of algebraic quantum gravity. While AQG is inspired by LQG, it differs drastically from it because in AQG there is fundamentally no topology or differential structure – it is background independent in a more generalized sense and could possibly have something to say about topology change. In this new formulation of quantum gravity AQG semiclassical states always control the fluctuations of all present degrees of freedom. This makes the AQG semiclassical analysis superior over that of LQG, and progress has been made in establishing it has the correct semiclassical limit and providing contact with familiar low energy physics. Physical applications of LQG Black hole entropy Black hole thermodynamics is the area of study that seeks to reconcile the laws of thermodynamics with the existence of black hole event horizons. The no hair conjecture of general relativity states that a black hole is characterized only by its mass, its charge, and its angular momentum; hence, it has no entropy. It appears, then, that one can violate the second law of thermodynamics by dropping an object with nonzero entropy into a black hole. Work by Stephen Hawking and Jacob Bekenstein showed that the second law of thermodynamics can be preserved by assigning to each black hole a black-hole entropy where is the area of the hole's event horizon, is the Boltzmann constant, and is the Planck length. The fact that the black hole entropy is also the maximal entropy that can be obtained by the Bekenstein bound (wherein the Bekenstein bound becomes an equality) was the main observation that led to the holographic principle. An oversight in the application of the no-hair theorem is the assumption that the relevant degrees of freedom accounting for the entropy of the black hole must be classical in nature; what if they were purely quantum mechanical instead and had non-zero entropy? This is what is realized in the LQG derivation of black hole entropy, and can be seen as a consequence of its background-independence – the classical black hole spacetime comes about from the semiclassical limit of the quantum state of the gravitational field, but there are many quantum states that have the same semiclassical limit. Specifically, in LQG it is possible to associate a quantum geometrical interpretation to the microstates: These are the quantum geometries of the horizon which are consistent with the area, , of the black hole and the topology of the horizon (i.e. spherical). LQG offers a geometric explanation of the finiteness of the entropy and of the proportionality of the area of the horizon. These calculations have been generalized to rotating black holes. It is possible to derive, from the covariant formulation of full quantum theory (Spinfoam) the correct relation between energy and area (1st law), the Unruh temperature and the distribution that yields Hawking entropy. The calculation makes use of the notion of dynamical horizon and is done for non-extremal black holes. A recent success of the theory in this direction is the computation of the entropy of all non singular black holes directly from theory and independent of Immirzi parameter. The result is the expected formula , where is the entropy and the area of the black hole, derived by Bekenstein and Hawking on heuristic grounds. This is the only known derivation of this formula from a fundamental theory, for the case of generic non singular black holes. Older attempts at this calculation had difficulties. The problem was that although Loop quantum gravity predicted that the entropy of a black hole is proportional to the area of the event horizon, the result depended on a crucial free parameter in the theory, the above-mentioned Immirzi parameter. However, there is no known computation of the Immirzi parameter, so it was fixed by demanding agreement with Bekenstein and Hawking's calculation of the black hole entropy. Hawking radiation in loop quantum gravity A detailed study of the quantum geometry of a black hole horizon has been made using loop quantum gravity. Loop-quantization does not reproduce the result for black hole entropy originally discovered by Bekenstein and Hawking, unless one chooses the value of the Immirzi parameter to cancel out another constant that arises in the derivation. However, it led to the computation of higher-order corrections to the entropy and radiation of black holes. Based on the fluctuations of the horizon area, a quantum black hole exhibits deviations from the Hawking spectrum that would be observable were X-rays from Hawking radiation of evaporating primordial black holes to be observed. The quantum effects are centered at a set of discrete and unblended frequencies highly pronounced on top of Hawking radiation spectrum. Planck star In 2014 Carlo Rovelli and Francesca Vidotto proposed that there is a Planck star inside every black hole. Based on LQG, the theory states that as stars are collapsing into black holes, the energy density reaches the Planck energy density, causing a repulsive force that creates a star. Furthermore, the existence of such a star would resolve the black hole firewall and black hole information paradox. Loop quantum cosmology The popular and technical literature makes extensive references to the LQG-related topic of loop quantum cosmology. LQC was mainly developed by Martin Bojowald. It was popularized in Scientific American for predicting a Big Bounce prior to the Big Bang. Loop quantum cosmology (LQC) is a symmetry-reduced model of classical general relativity quantized using methods that mimic those of loop quantum gravity (LQG) that predicts a "quantum bridge" between contracting and expanding cosmological branches. Achievements of LQC have been the resolution of the big bang singularity, the prediction of a Big Bounce, and a natural mechanism for inflation. LQC models share features of LQG and so is a useful toy model. However, the results obtained are subject to the usual restriction that a truncated classical theory, then quantized, might not display the true behaviour of the full theory due to artificial suppression of degrees of freedom that might have large quantum fluctuations in the full theory. It has been argued that singularity avoidance in LQC are by mechanisms only available in these restrictive models and that singularity avoidance in the full theory can still be obtained but by a more subtle feature of LQG. Loop quantum gravity phenomenology Quantum gravity effects are difficult to measure because the Planck length is so small. However recently physicists, such as Jack Palmer, have started to consider the possibility of measuring quantum gravity effects mostly from astrophysical observations and gravitational wave detectors. The energy of those fluctuations at scales this small cause space-perturbations which are visible at higher scales. Background-independent scattering amplitudes Loop quantum gravity is formulated in a background-independent language. No spacetime is assumed a priori, but rather it is built up by the states of theory themselves – however scattering amplitudes are derived from -point functions (Correlation function) and these, formulated in conventional quantum field theory, are functions of points of a background spacetime. The relation between the background-independent formalism and the conventional formalism of quantum field theory on a given spacetime is not obvious, and it is not obvious how to recover low-energy quantities from the full background-independent theory. One would like to derive the -point functions of the theory from the background-independent formalism, in order to compare them with the standard perturbative expansion of quantum general relativity and therefore check that loop quantum gravity yields the correct low-energy limit. A strategy for addressing this problem has been suggested; by studying the boundary amplitude, namely a path integral over a finite spacetime region, seen as a function of the boundary value of the field. In conventional quantum field theory, this boundary amplitude is well–defined and codes the physical information of the theory; it does so in quantum gravity as well, but in a fully background–independent manner. A generally covariant definition of -point functions can then be based on the idea that the distance between physical points – arguments of the -point function is determined by the state of the gravitational field on the boundary of the spacetime region considered. Progress has been made in calculating background-independent scattering amplitudes this way with the use of spin foams. This is a way to extract physical information from the theory. Claims to have reproduced the correct behaviour for graviton scattering amplitudes and to have recovered classical gravity have been made. "We have calculated Newton's law starting from a world with no space and no time." – Carlo Rovelli. Gravitons, string theory, supersymmetry, extra dimensions in LQG Some quantum theories of gravity posit a spin-2 quantum field that is quantized, giving rise to gravitons. In string theory, one generally starts with quantized excitations on top of a classically fixed background. This theory is thus described as background dependent. Particles like photons as well as changes in the spacetime geometry (gravitons) are both described as excitations on the string worldsheet. The background dependence of string theory can have physical consequences, such as determining the number of quark generations. In contrast, loop quantum gravity, like general relativity, is manifestly background independent, eliminating the background required in string theory. Loop quantum gravity, like string theory, also aims to overcome the nonrenormalizable divergences of quantum field theories. LQG does not introduce a background and excitations living on such a background, so LQG does not use gravitons as building blocks. Instead one expects that one may recover a kind of semiclassical limit or weak field limit where something like "gravitons" will show up again. In contrast, gravitons play a key role in string theory where they are among the first (massless) level of excitations of a superstring. LQG differs from string theory in that it is formulated in 3 and 4 dimensions and without supersymmetry or Kaluza–Klein extra dimensions, while the latter requires both to be true. There is no experimental evidence to date that confirms string theory's predictions of supersymmetry and Kaluza–Klein extra dimensions. In a 2003 paper "A Dialog on Quantum Gravity", Carlo Rovelli regards the fact LQG is formulated in 4 dimensions and without supersymmetry as a strength of the theory as it represents the most parsimonious explanation, consistent with current experimental results, over its rival string/M-theory. Proponents of string theory will often point to the fact that, among other things, it demonstrably reproduces the established theories of general relativity and quantum field theory in the appropriate limits, which loop quantum gravity has struggled to do. In that sense string theory's connection to established physics may be considered more reliable and less speculative, at the mathematical level. Loop quantum gravity has nothing to say about the matter (fermions) in the universe. Since LQG has been formulated in 4 dimensions (with and without supersymmetry), and M-theory requires supersymmetry and 11 dimensions, a direct comparison between the two has not been possible. It is possible to extend mainstream LQG formalism to higher-dimensional supergravity, general relativity with supersymmetry and Kaluza–Klein extra dimensions should experimental evidence establish their existence. It would therefore be desirable to have higher-dimensional Supergravity loop quantizations at one's disposal in order to compare these approaches. A series of papers have been published attempting this. Most recently, Thiemann (and alumni) have made progress toward calculating black hole entropy for supergravity in higher dimensions. It will be useful to compare these results to the corresponding super string calculations. LQG and related research programs Several research groups have attempted to combine LQG with other research programs: Johannes Aastrup, Jesper M. Grimstrup et al. research combines noncommutative geometry with canonical quantum gravity and Ashtekar variables, Laurent Freidel, Simone Speziale, et al., spinors and twistor theory with loop quantum gravity, and Lee Smolin et al. with Verlinde entropic gravity and loop gravity. Stephon Alexander, Antonino Marciano and Lee Smolin have attempted to explain the origins of weak force chirality in terms of Ashketar's variables, which describe gravity as chiral, and LQG with Yang–Mills theory fields in four dimensions. Sundance Bilson-Thompson, Hackett et al., has attempted to introduce the standard model via LQGs degrees of freedom as an emergent property (by employing the idea of noiseless subsystems, a notion introduced in a more general situation for constrained systems by Fotini Markopoulou-Kalamara et al.) Furthermore, LQG has drawn philosophical comparisons with causal dynamical triangulation and asymptotically safe gravity, and the spinfoam with group field theory and AdS/CFT correspondence. Smolin and Wen have suggested combining LQG with string-net liquid, tensors, and Smolin and Fotini Markopoulou-Kalamara quantum graphity. There is the consistent discretizations approach. Also, Pullin and Gambini provide a framework to connect the path integral and canonical approaches to quantum gravity. They may help reconcile the spin foam and canonical loop representation approaches. Recent research by Chris Duston and Matilde Marcolli introduces topology change via topspin networks. Problems and comparisons with alternative approaches Some of the major unsolved problems in physics are theoretical, meaning that existing theories seem incapable of explaining a certain observed phenomenon or experimental result. The others are experimental, meaning that there is a difficulty in creating an experiment to test a proposed theory or investigate a phenomenon in greater detail. Many of these problems apply to LQG, including: Can quantum mechanics and general relativity be realized as a fully consistent theory (perhaps as a quantum field theory)? Is spacetime fundamentally continuous or discrete? Would a consistent theory involve a force mediated by a hypothetical graviton, or be a product of a discrete structure of spacetime itself (as in loop quantum gravity)? Are there deviations from the predictions of general relativity at very small or very large scales or in other extreme circumstances that flow from a quantum gravity theory? The theory of LQG is one possible solution to the problem of quantum gravity, as is string theory. There are substantial differences however. For example, string theory also addresses unification, the understanding of all known forces and particles as manifestations of a single entity, by postulating extra dimensions and so-far unobserved additional particles and symmetries. Contrary to this, LQG is based only on quantum theory and general relativity and its scope is limited to understanding the quantum aspects of the gravitational interaction. On the other hand, the consequences of LQG are radical, because they fundamentally change the nature of space and time and provide a tentative but detailed physical and mathematical picture of quantum spacetime. Presently, no semiclassical limit recovering general relativity has been shown to exist. This means it remains unproven that LQG's description of spacetime at the Planck scale has the right continuum limit (described by general relativity with possible quantum corrections). Specifically, the dynamics of the theory are encoded in the Hamiltonian constraint, but there is no candidate Hamiltonian. Other technical problems include finding off-shell closure of the constraint algebra and physical inner product vector space, coupling to matter fields of quantum field theory, fate of the renormalization of the graviton in perturbation theory that lead to ultraviolet divergence beyond 2-loops (see one-loop Feynman diagram in Feynman diagram). While there has been a proposal relating to observation of naked singularities, and doubly special relativity as a part of a program called loop quantum cosmology, there is no experimental observation for which loop quantum gravity makes a prediction not made by the Standard Model or general relativity (a problem that plagues all current theories of quantum gravity). Because of the above-mentioned lack of a semiclassical limit, LQG has not yet even reproduced the predictions made by general relativity. An alternative criticism is that general relativity may be an effective field theory, and therefore quantization ignores the fundamental degrees of freedom. ESA's INTEGRAL satellite measured polarization of photons of different wavelengths and was able to place a limit in the granularity of space that is less than 10−48m or 13 orders of magnitude below the Planck scale.	Physical sciences	Quantum mechanics	Physics
152671	https://en.wikipedia.org/wiki/Z3%20%28computer%29	Z3 (computer)	The Z3 was a German electromechanical computer designed by Konrad Zuse in 1938, and completed in 1941. It was the world's first working programmable, fully automatic digital computer. The Z3 was built with 2,600 relays, implementing a 22-bit word length that operated at a clock frequency of about 5–10 Hz. Program code was stored on punched film. Initial values were entered manually. The Z3 was completed in Berlin in 1941. It was not considered vital, so it was never put into everyday operation. Based on the work of the German aerodynamics engineer Hans Georg Küssner (known for the Küssner effect), a "Program to Compute a Complex Matrix" was written and used to solve wing flutter problems. Zuse asked the German government for funding to replace the relays with fully electronic switches, but funding was denied during World War II since such development was deemed "not war-important". The original Z3 was destroyed on 21 December 1943 during an Allied bombardment of Berlin. That Z3 was originally called V3 (Versuchsmodell 3 or Experimental Model 3) but was renamed so that it would not be confused with Germany's V-weapons. A fully functioning replica was built in 1961 by Zuse's company, Zuse KG, which is now on permanent display at Deutsches Museum in Munich. The Z3 was demonstrated in 1998 to be, in principle, Turing-complete. However, because it lacked conditional branching, the Z3 only meets this definition by speculatively computing all possible outcomes of a calculation. Thanks to this machine and its predecessors, Konrad Zuse has often been suggested as the inventor of the computer. Design and development Zuse designed the Z1 in 1935 to 1936 and built it from 1936 to 1938. The Z1 was wholly mechanical and only worked for a few minutes at a time at most. Helmut Schreyer advised Zuse to use a different technology. As a doctoral student at the Technische Hochschule in Charlottenburg (now Technische Universität Berlin) in 1937 he worked on the implementation of Boolean operations and (in today's terminology) flip-flops on the basis of vacuum tubes. In 1938, Schreyer demonstrated a circuit on this basis to a small audience, and explained his vision of an electronic computing machine – but since the largest operational electronic devices contained far fewer tubes this was considered practically infeasible. In that year when presenting the plan for a computer with 2,000 electron tubes, Zuse and Schreyer, who was an assistant at Telecommunication Institute at Technische Universität Berlin, were discouraged by members of the institute who knew about the problems with electron tube technology. Zuse later recalled: "They smiled at us in 1939, when we wanted to build electronic machines ... We said: The electronic machine is great, but first the components have to be developed." In 1940, Zuse and Schreyer managed to arrange a meeting at the Oberkommando der Wehrmacht (OKW) to discuss a potential project for developing an electronic computer, but when they estimated a duration of two or three years, the proposal was rejected. Zuse decided to implement the next design based on relays. The realization of the Z2 was helped financially by Kurt Pannke, who manufactured small calculating machines. The Z2 was completed and presented to an audience of the ("German Laboratory for Aviation") in 1940 in Berlin-Adlershof. Zuse was lucky – this presentation was one of the few instances where the Z2 actually worked and could convince the DVL to partly finance the next design. In 1941, improving on the basic Z2 machine, he built the Z3 in a highly secret project of the German government. Joseph Jennissen (1905–1977), member of the "Research-Leadership" (Forschungsführung) in the Reich Air Ministry acted as a government supervisor for orders of the ministry to Zuse's company ZUSE Apparatebau. A further intermediary between Zuse and the Reich Air Ministry was the aerodynamicist Herbert A. Wagner. The Z3 was completed in 1941 and was faster and far more reliable than the Z1 and Z2. The Z3 floating-point arithmetic was improved over that of the Z1 in that it implemented exception handling "using just a few relays", the exceptional values (plus infinity, minus infinity and undefined) could be generated and passed through operations. It further added a square root instruction. The Z3, like its predecessors, stored its program on an external punched tape, thus no rewiring was necessary to change programs. However, it did not have conditional branching found in later universal computers. On 12 May 1941, the Z3 was presented to an audience of scientists including the professors Alfred Teichmann and Curt Schmieden of the ("German Laboratory for Aviation") in Berlin, today known as the German Aerospace Center in Cologne. Zuse moved on to the Z4 design, which he completed in a bunker in the Harz mountains, alongside Wernher von Braun's ballistic missile development. When World War II ended, Zuse retreated to Hinterstein in the Alps with the Z4, where he remained for several years. Instruction set The Z3 operated as a stack machine with a stack of two registers, R1 and R2. The first load operation in a program would load the contents of a memory location into R1; the next load operation would load the contents of a memory location into R2. Arithmetic instructions would operate on the contents of R1 and R2, leaving the result in R1, and clearing R2; the next load operation would load into R2. A store operation would store the contents of R1 into a memory location, and clear R1; the next load operation would load the contents of a memory location into R1. A read keyboard operation would read a number from the keyboard into R1 and clear R2. A display instruction would display the contents of R1 and clear R2; the next load instruction would load into R2. Z3 as a universal Turing machine It was possible to construct loops on the Z3, but there was no conditional branch instruction. Nevertheless, the Z3 was Turing-complete – how to implement a universal Turing machine on the Z3 was shown in 1998 by Raúl Rojas. He proposed that the tape program would have to be long enough to execute every possible path through both sides of every branch. It would compute all possible answers, but the unneeded results would be canceled out (a kind of speculative execution). Rojas concludes, "We can therefore say that, from an abstract theoretical perspective, the computing model of the Z3 is equivalent to the computing model of today's computers. From a practical perspective, and in the way the Z3 was really programmed, it was not equivalent to modern computers." This seeming limitation belies the fact that the Z3 provided a practical instruction set for the typical engineering applications of the 1940s. Mindful of the existing hardware restrictions, Zuse's main goal at the time was to have a workable device to facilitate his work as a civil engineer. Relation to other work The success of Zuse's Z3 is often attributed to its use of the simple binary system. This was invented roughly three centuries earlier by Gottfried Leibniz; Boole later used it to develop his Boolean algebra. Zuse was inspired by Hilbert's and Ackermann's book on elementary mathematical logic Principles of Mathematical Logic. In 1937, Claude Shannon introduced the idea of mapping Boolean algebra onto electronic relays in a seminal work on digital circuit design. Zuse, however, did not know of Shannon's work and developed the groundwork independently for his first computer Z1, which he designed and built from 1935 to 1938. Zuse's coworker Helmut Schreyer built an electronic digital experimental model of a computer using 100 vacuum tubes in 1942, but it was lost at the end of the war. An analog computer was built by the rocket scientist Helmut Hölzer in 1942 at the Peenemünde Army Research Center to simulate V-2 rocket trajectories. The Colossus (1943), built by Tommy Flowers, and the Atanasoff–Berry computer (1942) used thermionic valves (vacuum tubes) and binary representation of numbers. Programming was by means of re-plugging patch panels and setting switches. The ENIAC computer, completed after the war, used vacuum tubes to implement switches and used decimal representation for numbers. Until 1948 programming was, as with Colossus, by patch leads and switches. The Manchester Baby of 1948 along with the Manchester Mark 1 and EDSAC both of 1949 were the world's earliest working computers that stored program instructions and data in the same space. In this they implemented the stored-program concept which is frequently (but erroneously) attributed to a 1945 paper by John von Neumann and colleagues. Von Neumann is said to have given due credit to Alan Turing, and the concept had actually been mentioned earlier by Konrad Zuse himself, in a 1936 patent application (that was rejected). Konrad Zuse himself remembered in his memoirs: "During the war it would have barely been possible to build efficient stored program devices anyway." Friedrich L. Bauer later wrote: "His visionary ideas (live programs) which were only to be published years afterwards aimed at the right practical direction but were never implemented by him." Specifications Average calculation speed: addition – 0.8 seconds, multiplication – 3 seconds Arithmetic unit: Binary floating-point, 22-bit, add, subtract, multiply, divide, square root Data memory: 64 22-bit words Program memory: Punched celluloid tape Input: Decimal floating-point numbers Output: Decimal floating-point numbers Input and Output was facilitated by a terminal, with a special keyboard for input and a row of lamps to show results Elements: Around 2,000 relays (1,400 for the memory) Frequency: 5–10 hertz Power consumption: Around 4,000 watts Weight: Around Modern reconstructions A modern reconstruction directed by Raúl Rojas and Horst Zuse started in 1997 and finished in 2003. It is now in the Konrad Zuse Museum in Hünfeld, Germany. Memory was halved to 32 words. Power consumption is about 400 W, and weight is about . In 2008, Horst Zuse started a reconstruction of the Z3 by himself. It was presented in 2010 in the Konrad Zuse Museum in Hünfeld.	Technology	Early computers	null
152690	https://en.wikipedia.org/wiki/Hovercraft	Hovercraft	A hovercraft (: hovercraft), also known as an air-cushion vehicle or ACV, is an amphibious craft capable of travelling over land, water, mud, ice, and various other surfaces. Hovercraft use blowers to produce a large volume of air below the hull, or air cushion, that is slightly above atmospheric pressure. The pressure difference between the higher-pressure air below the hull and lower pressure ambient air above it produces lift, which causes the hull to float above the running surface. For stability reasons, the air is typically blown through slots or holes around the outside of a disk- or oval-shaped platform, giving most hovercraft a characteristic rounded-rectangle shape. The first practical design for hovercraft was derived from a British invention in the 1950s. They are now used throughout the world as specialised transports in disaster relief, coastguard, military and survey applications, as well as for sport or passenger service. Very large versions have been used to transport hundreds of people and vehicles across the English Channel, whilst others have military applications used to transport tanks, soldiers and large equipment in hostile environments and terrain. Decline in public demand meant that , the only year-round public hovercraft service in the world still in operation serves between the Isle of Wight and Southsea in the UK. Oita Hovercraft is planning to resume services in Oita, Japan in 2024. Although now a generic term for the type of craft, the name Hovercraft itself was a trademark owned by Saunders-Roe (later British Hovercraft Corporation (BHC), then Westland), hence other manufacturers' use of alternative names to describe the vehicles. History Early efforts There have been many attempts to understand the principles of high air pressure below hulls and wings. Hovercraft are unique in that they can lift themselves while still, differing from ground effect vehicles and hydrofoils that require forward motion to create lift. The first mention, in the historical record of the concepts behind surface-effect vehicles, to use the term hovering was by Swedish scientist Emanuel Swedenborg in 1716. The shipbuilder John Isaac Thornycroft patented an early design for an air cushion ship / hovercraft in the 1870s, but suitable, powerful, engines were not available until the 20th century. In 1915, the Austrian Dagobert Müller von Thomamühl (1880–1956) built the world's first "air cushion" boat (). Shaped like a section of a large aerofoil (this creates a low-pressure area above the wing much like an aircraft), the craft was propelled by four aero engines driving two submerged marine propellers, with a fifth engine that blew air under the front of the craft to increase the air pressure under it. Only when in motion could the craft trap air under the front, increasing lift. The vessel also required a depth of water to operate and could not transition to land or other surfaces. Designed as a fast torpedo boat, the had a top speed of over . It was thoroughly tested and even armed with torpedoes and machine guns for operation in the Adriatic. It never saw actual combat, however, and as the war progressed it was eventually scrapped due to a lack of interest and perceived need, and its engines returned to the air force. The theoretical grounds for motion over an air layer were constructed by Konstantin Eduardovich Tsiolkovskii in 1926 and 1927. In 1929, Andrew Kucher of Ford began experimenting with the Levapad concept, metal disks with pressurized air blown through a hole in the centre. Levapads do not offer stability on their own. Several must be used together to support a load above them. Lacking a skirt, the pads had to remain very close to the running surface. He initially imagined these being used in place of casters and wheels in factories and warehouses, where the concrete floors offered the smoothness required for operation. By the 1950s, Ford showed a number of toy models of cars using the system, but mainly proposed its use as a replacement for wheels on trains, with the Levapads running close to the surface of existing rails. In 1931, Finnish aero engineer Toivo J. Kaario began designing a developed version of a vessel using an air cushion and built a prototype ('Surface Glider'), in 1937. His design included the modern features of a lift engine blowing air into a flexible envelope for lift. Kaario's efforts were followed closely in the Soviet Union by Vladimir Levkov, who returned to the solid-sided design of the . Levkov designed and built a number of similar craft during the 1930s, and his L-5 fast-attack boat reached in testing. However, the start of World War II put an end to his development work. During World War II, an American engineer, Charles Fletcher, invented a walled air cushion vehicle, the Glidemobile. Because the project was classified by the U.S. government, Fletcher could not file a patent. In April 1958, Ford engineers demonstrated the Glide-air, a model of a wheel-less vehicle that speeds on a thin film of air only 76.2 μm ( of an inch) above its tabletop roadbed. An article in Modern Mechanix quoted Andrew A. Kucher, Ford's vice president in charge of Engineering and Research noting "We look upon Glide-air as a new form of high-speed land transportation, probably in the field of rail surface travel, for fast trips of distances of up to about ". In 1959, Ford displayed a hovercraft concept car, the Ford Levacar Mach I. In August 1961, Popular Science reported on the Aeromobile 35B, an air-cushion vehicle (ACV) that was invented by William R. Bertelsen and was envisioned to revolutionise the transportation system, with personal hovering self-driving cars that could speed up to . Christopher Cockerell The idea of the modern hovercraft is most often associated with Christopher Cockerell, a British mechanical engineer. Cockerell's group was the first to develop the use of a ring of air for maintaining the cushion, the first to develop a successful skirt, and the first to demonstrate a practical vehicle in continued use. A memorial to Cockerell's first design stands in the village of Somerleyton. Cockerell came across the key concept in his design when studying the ring of airflow when high-pressure air was blown into the annular area between two concentric tin cans (one coffee and the other from cat food) and a hairdryer. This produced a ring of airflow, as expected, but he noticed an unexpected benefit as well; the sheet of fast-moving air presented a sort of physical barrier to the air on either side of it. This effect, which he called the "momentum curtain", could be used to trap high-pressure air in the area inside the curtain, producing a high-pressure plenum that earlier examples had to build up with considerably more airflow. In theory, only a small amount of active airflow would be needed to create lift and much less than a design that relied only on the momentum of the air to provide lift, like a helicopter. In terms of power, a hovercraft would only need between one quarter to one half of the power required by a helicopter. Cockerell built and tested several models of his hovercraft design in Somerleyton, Suffolk, during the early 1950s. The design featured an engine mounted to blow from the front of the craft into a space below it, combining both lift and propulsion. He demonstrated the model flying over many Whitehall carpets in front of various government experts and ministers, and the design was subsequently put on the secret list. In spite of tireless efforts to arrange funding, no branch of the military was interested, as he later joked, "The Navy said it was a plane not a boat; the RAF said it was a boat not a plane; and the Army were 'plain not interested'." SR.N1 This lack of military interest meant that there was no reason to keep the concept secret, and it was declassified. Cockerell was finally able to convince the National Research Development Corporation to fund development of a full-scale model. In 1958, the NRDC placed a contract with Saunders-Roe for the development of what would become the SR.N1, short for "Saunders-Roe, Nautical 1". The SR.N1 was powered by a 450 hp Alvis Leonides engine powering a vertical fan in the middle of the craft. In addition to providing the lift air, a portion of the airflow was bled off into two channels on either side of the craft, which could be directed to provide thrust. In normal operation this extra airflow was directed rearward for forward thrust and blew over two large vertical rudders that provided directional control. For low-speed manoeuvrability, the extra thrust could be directed fore or aft, differentially for rotation. The SR.N1 made its first hover on 11 June 1959, and made its famed successful crossing of the English Channel on 25 July 1959. In December 1959, the Duke of Edinburgh visited Saunders-Roe at East Cowes and persuaded the chief test-pilot, Commander Peter Lamb, to allow him to take over the SR.N1's controls. He flew the SR.N1 so fast that he was asked to slow down a little. On examination of the craft afterwards, it was found that she had been dished in the bow due to excessive speed, damage that was never allowed to be repaired, and was from then on affectionately referred to as the 'Royal Dent'. Skirts and other improvements Testing quickly demonstrated that the idea of using a single engine to provide air for both the lift curtain and forward flight required too many trade-offs. A Blackburn Marboré turbojet for forward thrust and two large vertical rudders for directional control were added, producing the SR.N1 Mk II. A further upgrade with the Armstrong Siddeley Viper produced the Mk III. Further modifications, especially the addition of pointed nose and stern areas, produced the Mk IV. Although the SR.N1 was successful as a testbed, the design hovered too close to the surface to be practical; at even small waves would hit the bow. The solution was offered by Cecil Latimer-Needham, following a suggestion made by his business partner Arthur Ord-Hume. In 1958, he suggested the use of two rings of rubber to produce a double-walled extension of the vents in the lower fuselage. When air was blown into the space between the sheets it exited the bottom of the skirt in the same way it formerly exited the bottom of the fuselage, re-creating the same momentum curtain, but this time at some distance from the bottom of the craft. Latimer-Needham and Cockerell devised a high skirt design, which was fitted to the SR.N1 to produce the Mk V, displaying hugely improved performance, with the ability to climb over obstacles almost as high as the skirt. In October 1961, Latimer-Needham sold his skirt patents to Westland, who had recently taken over Saunders Roe's interest in the hovercraft. Experiments with the skirt design demonstrated a problem; it was originally expected that pressure applied to the outside of the skirt would bend it inward, and the now-displaced airflow would cause it to pop back out. What actually happened is that the slight narrowing of the distance between the walls resulted in less airflow, which in turn led to more air loss under that section of the skirt. The fuselage above this area would drop due to the loss of lift at that point, and this led to further pressure on the skirt. After considerable experimentation, Denys Bliss at Hovercraft Development Ltd. found the solution to this problem. Instead of using two separate rubber sheets to form the skirt, a single sheet of rubber was bent into a U shape to provide both sides, with slots cut into the bottom of the U forming the annular vent. When deforming pressure was applied to the outside of this design, air pressure in the rest of the skirt forced the inner wall to move in as well, keeping the channel open. Although there was some deformation of the curtain, the airflow within the skirt was maintained and the lift remained relatively steady. Over time, this design evolved into individual extensions over the bottom of the slots in the skirt, known as "fingers". Commercialisation Through these improvements, the hovercraft became an effective transport system for high-speed service on water and land, leading to widespread developments for military vehicles, search and rescue, and commercial operations. By 1962, many UK aviation and shipbuilding firms were working on hovercraft designs, including Saunders Roe/Westland, Vickers-Armstrong, William Denny, Britten-Norman and Folland. Small-scale ferry service started as early as 1962 with the launch of the Vickers-Armstrong VA-3. With the introduction of the 254 passenger and 30 car carrying SR.N4 cross-channel ferry by Hoverlloyd and Seaspeed in 1968, hovercraft had developed into useful commercial craft. Another major pioneering effort of the early hovercraft era was carried out by Jean Bertin's firm in France. Bertin was an advocate of the "multi-skirt" approach, which used a number of smaller cylindrical skirts instead of one large one in order to avoid the problems noted above. During the early 1960s he developed a series of prototype designs, which he called "terraplanes" if they were aimed for land use, and "naviplanes" for water. The best known of these designs was the N500 Naviplane, built for Seaspeed by the Société d'Etude et de Développement des Aéroglisseurs Marins (SEDAM). The N500 could carry 400 passengers, 55 cars and five buses. It set a speed record between Boulogne and Dover of . It was rejected by its operators, who claimed that it was unreliable. Another discovery was that the total amount of air needed to lift the craft was a function of the roughness of the surface over which it travelled. On flat surfaces, like pavement, the required air pressure was so low that hovercraft were able to compete in energy terms with conventional systems like steel wheels. However, the hovercraft lift system acted as both a lift and a very effective suspension, and thus it naturally lent itself to high-speed use where conventional suspension systems were considered too complex. This led to a variety of "hovertrain" proposals during the 1960s, including England's Tracked Hovercraft and France's Aérotrain. In the U.S., Rohr Inc. and Garrett both took out licences to develop local versions of the Aérotrain. These designs competed with maglev systems in the high-speed arena, where their primary advantage was the very "low tech" tracks they needed. On the downside, the air blowing dirt and trash out from under the trains presented a unique problem in stations, and interest in them waned in the 1970s. By the early 1970s, the basic concept had been well developed, and the hovercraft had found a number of niche roles where its combination of features were advantageous. Today, they are found primarily in military use for amphibious operations, search-and-rescue vehicles in shallow water, and sporting vehicles. Design Hovercraft can be powered by one or more engines. Smaller craft, such as the SR.N6, usually have one engine with the drive split through a gearbox. On vehicles with several engines, one usually drives the fan (or impeller), which is responsible for lifting the vehicle by forcing high pressure air under the craft. The air inflates the "skirt" under the vehicle, causing it to rise above the surface. Additional engines provide thrust in order to propel the craft. Some hovercraft use ducting to allow one engine to perform both tasks by directing some of the air to the skirt, the rest of the air passing out of the back to push the craft forward. Uses Commercial The British aircraft and marine engineering company Saunders-Roe built the first practical human-carrying hovercraft for the National Research Development Corporation, the SR.N1, which carried out several test programmes in 1959 to 1961 (the first public demonstration was in 1959), including a cross-channel test run in July 1959, piloted by Peter "Sheepy" Lamb, an ex-naval test pilot and the chief test pilot at Saunders Roe. Christopher Cockerell was on board, and the flight took place on the 50th anniversary of Louis Blériot's first aerial crossing. The SR.N1 was driven by expelled air, powered by a single piston engine. Demonstrated at the Farnborough Airshow in 1960, it was shown that this simple craft can carry a load of up to 12 marines with their equipment as well as the pilot and co-pilot with only a slight reduction in hover height proportional to the load carried. The SR.N1 did not have any skirt, using instead the peripheral air principle that Cockerell had patented. It was later found that the craft's hover height was improved by the addition of a skirt of flexible fabric or rubber around the hovering surface to contain the air. The skirt was an independent invention made by a Royal Navy officer, C.H. Latimer-Needham, who sold his idea to Westland (by then the parent of Saunders-Roe's helicopter and hovercraft interests), and who worked with Cockerell to develop the idea further. The first passenger-carrying hovercraft to enter service was the Vickers VA-3, which, in the summer of 1962, carried passengers regularly along the north Wales coast from Moreton, Merseyside, to Rhyl. It was powered by two turboprop aero-engines and driven by propellers. During the 1960s, Saunders-Roe developed several larger designs that could carry passengers, including the SR.N2, which operated across the Solent, in 1962, and later the SR.N6, which operated across the Solent from Southsea to Ryde on the Isle of Wight for many years. In 1963 the SR.N2 was used in experimental service between Weston-super-Mare and Penarth under the aegis of P & A Campbell, the paddle steamer operators. Operations by Hovertravel commenced on 24 July 1965, using the SR.N6, which carried 38 passengers. Two 98 seat AP1-88 hovercraft were introduced on this route in 1983, and in 2007, these were joined by the first 130-seat BHT130 craft. The AP1-88 and the BHT130 were notable as they were largely built by Hoverwork using shipbuilding techniques and materials (i.e. welded aluminium structure and diesel engines) rather than the aircraft techniques used to build the earlier craft built by Saunders-Roe-British Hovercraft Corporation. Over 20 million passengers had used the service as of 2004 – the service is still operating () and is by far the longest, continuously-operated hovercraft service. In 1966, two cross-channel passenger hovercraft services were inaugurated using SR.N6 hovercraft. Hoverlloyd ran services from Ramsgate Harbour, England, to Calais, France, and Townsend Ferries also started a service to Calais from Dover, which was soon superseded by that of Seaspeed. As well as Saunders-Roe and Vickers (which combined in 1966 to form the British Hovercraft Corporation (BHC)), other commercial craft were developed during the 1960s in the UK by Cushioncraft (part of the Britten-Norman Group) and Hovermarine based at Woolston (the latter being sidewall hovercraft, where the sides of the hull projected down into the water to trap the cushion of air with normal hovercraft skirts at the bow and stern). One of these models, the HM-2, was used by Red Funnel between Southampton (near the Woolston Floating Bridge) and Cowes. The world's first car-carrying hovercraft was made in 1968, the BHC Mountbatten class (SR.N4) models, each powered by four Bristol Proteus turboshaft engines. These were both used by rival operators Hoverlloyd and Seaspeed (which joined to form Hoverspeed in 1981) to operate regular car and passenger carrying services across the English Channel. Hoverlloyd operated from Ramsgate, where a special hoverport had been built at Pegwell Bay, to Calais. Seaspeed operated from Dover, England, to Calais and Boulogne in France. The first SR.N4 had a capacity of 254 passengers and 30 cars, and a top speed of . The channel crossing took around 30 minutes and was run like an airline with flight numbers. The later SR.N4 Mk.III had a capacity of 418 passengers and 60 cars. These were later joined by the French-built SEDAM N500 Naviplane with a capacity of 385 passengers and 45 cars; only one entered service and was used intermittently for a few years on the cross-channel service until returned to SNCF in 1983. The service ceased on 1 October 2000 after 32 years, due to competition with traditional ferries, catamarans, the disappearance of duty-free shopping within the EU, the advancing age of the SR.N4 hovercraft, and the opening of the Channel Tunnel. The commercial success of hovercraft suffered from rapid rises in fuel prices during the late 1960s and 1970s, following conflict in the Middle East. Alternative over-water vehicles, such as wave-piercing catamarans (marketed as the SeaCat in the UK until 2005), use less fuel and can perform most of the hovercraft's marine tasks. Although developed elsewhere in the world for both civil and military purposes, except for the Solent Ryde-to-Southsea crossing, hovercraft disappeared from the coastline of Britain until a range of Griffon Hoverwork were bought by the Royal National Lifeboat Institution. Hovercraft used to ply between the Gateway of India in Mumbai and CBD Belapur and Vashi in Navi Mumbai between 1994 and 1999, but the services were subsequently stopped due to the lack of sufficient water transport infrastructure. Civilian non-commercial In Finland, small hovercraft are widely used in maritime rescue and during the rasputitsa ("mud season") as archipelago liaison vehicles. In England, hovercraft of the Burnham-on-Sea Area Rescue Boat (BARB) are used to rescue people from thick mud in Bridgwater Bay. Avon Fire and Rescue Service became the first Local Authority fire service in the UK to operate a hovercraft. It is used to rescue people from thick mud in the Weston-super-Mare area and during times of inland flooding. A Griffon rescue hovercraft has been in use for a number of years with the Airport Fire Service at Dundee Airport in Scotland. It is used in the event of an aircraft ditching in the Tay estuary. Numerous fire departments around the US/Canadian Great Lakes operate hovercraft for water and ice rescues, often of ice fisherman stranded when ice breaks off from shore. The Canadian Coast Guard uses hovercraft to break light ice. In October 2008, The Red Cross commenced a flood-rescue service hovercraft based in Inverness, Scotland. Gloucestershire Fire and Rescue Service received two flood-rescue hovercraft donated by Severn Trent Water following the 2007 UK floods. Since 2006, hovercraft have been used in aid in Madagascar by HoverAid, an international NGO who use the hovercraft to reach the most remote places on the island. The Scandinavian airline SAS used to charter an AP1-88 hovercraft for regular passengers between Copenhagen Airport, Denmark, and the SAS Hovercraft Terminal in Malmö, Sweden. In 1998, the US Postal Service began using the British built Hoverwork AP1-88 to haul mail, freight, and passengers from Bethel, Alaska, to and from eight small villages along the Kuskokwim River. Bethel is far removed from the Alaska road system, thus making the hovercraft an attractive alternative to the air based delivery methods used prior to introduction of the hovercraft service. Hovercraft service is suspended for several weeks each year while the river is beginning to freeze to minimize damage to the river ice surface. The hovercraft is able to operate during the freeze-up period; however, this could potentially break the ice and create hazards for villagers using their snowmobiles along the river during the early winter. In 2006, Kvichak Marine Industries of Seattle, US built, under licence, a cargo/passenger version of the Hoverwork BHT130. Designated 'Suna-X', it is used as a high-speed ferry for up to 47 passengers and of freight serving the remote Alaskan villages of King Cove and Cold Bay. An experimental service was operated in Scotland across the Firth of Forth (between Kirkcaldy and Portobello, Edinburgh), from 16 to 28 July 2007. Marketed as Forthfast, the service used a craft chartered from Hovertravel and achieved an 85% passenger load factor. , the possibility of establishing a permanent service is still under consideration. Since the Channel routes abandoned hovercraft, and pending any reintroduction on the Scottish route, the United Kingdom's only public hovercraft service is that operated by Hovertravel between Southsea (Portsmouth) and Ryde on the Isle of Wight. From the 1960s, several commercial lines were operated in Japan, without much success. In Japan the last commercial line had linked Ōita Airport and central Ōita but was shut down in October 2009. However, the commercial line between Ōita Airport and central Ōita is scheduled to reopen in 2024. Hovercraft are still manufactured in the UK, near to where they were first conceived and tested, on the Isle of Wight. They can also be chartered for a wide variety of uses including inspections of shallow bed offshore wind farms and VIP or passenger use. A typical vessel would be a Tiger IV or a Griffon. They are light, fast, road transportable and very adaptable with the unique feature of minimising damage to environments. Military China The People's Army Navy of China operates the Jingsah II class LCAC. This troop and equipment carrying hovercraft is roughly the Chinese equivalent of the U.S. Navy LCAC. Finland The Finnish Navy designed an experimental missile attack hovercraft class, Tuuli class hovercraft, in the late 1990s. The prototype of the class, Tuuli, was commissioned in 2000. It proved an extremely successful design for a littoral fast attack craft, but due to fiscal reasons and doctrinal change in the Navy, the hovercraft was soon withdrawn. Iran The Iranian Navy operates multiple British-made and some Iranian-produced hovercraft. The Tondar or Thunderbolt comes in varieties designed for combat and transportation. Iran has equipped the Tondar with mid-range missiles, machine guns and retrievable reconnaissance drones. Currently they are used for water patrols and combat against drug smugglers. Russia The Soviet Union had military hovercraft such as the small Czilim-class hovercraft, comparable to the SR.N6, the large Zubr-class landing craft and the Bora missile launcher surface effect ship (a hybrid between a hovercraft and a catamaran). United Kingdom The first application of the hovercraft for military use was by the British Armed Forces, using hovercraft built by Saunders-Roe. In 1961, the United Kingdom set up the Interservice Hovercraft Trials Unit (IHTU) based at RNAS Lee-on-Solent (HMS Daedalus), now the site of the Hovercraft Museum, near Portsmouth. This unit carried out trials on the SR.N1 from Mk1 through Mk5 as well as testing the SR.N2, SR.N3, SR.N5 and SR.N6 craft. The Hovercraft Trials Unit (Far East) was established by the Royal Navy at Singapore in August 1964 with two armed hovercraft; they were deployed later that year to Tawau in Malaysian Borneo and operated on waterways there during the Indonesia–Malaysia confrontation. The hovercraft's inventor, Sir Christopher Cockerell, claimed late in his life that the Falklands War could have been won far more easily had the British military shown more commitment to the hovercraft; although earlier trials had been conducted in the Falkland Islands with an SRN-6, the hovercraft unit had been disbanded by the time of the conflict. Currently, the Royal Marines use the Griffonhoverwork 2400TD hovercraft, the replacement for the Griffon 2000 TDX Class ACV, which was deployed operationally by the marines in the 2003 invasion of Iraq. United States During the 1960s, Bell licensed and sold the Saunders-Roe SR.N5 as the Bell SK-5. They were deployed on trial to the Vietnam War by the United States Navy as PACV patrol craft in the Mekong Delta where their mobility and speed was unique. This was used in both the UK SR.N5 curved deck configuration and later with modified flat deck, gun turret and grenade launcher designated the 9255 PACV. The United States Army also experimented with the use of SR.N5 hovercraft in Vietnam. Three hovercraft with the flat deck configuration were deployed to Đồng Tâm in the Mekong Delta region and later to Ben Luc. They saw action primarily in the Plain of Reeds. One was destroyed in early 1970 and another in August of that same year, after which the unit was disbanded. The only remaining U.S. Army SR.N5 hovercraft is currently on display in the Army Transport Museum in Virginia. Experience led to the proposed Bell SK-10, which was the basis for the LCAC-class air-cushioned landing craft now deployed by the U.S. and Japanese Navy. Developed and tested in the mid-1970s, the LACV-30 was used by the US Army to transport military cargo in logistics-over-the-shore operations from the early 1980s until the mid-1990s. Recreational/sport Small commercially manufactured, kit or plan-built hovercraft are increasingly being used for recreational purposes, such as inland racing and cruising on inland lakes and rivers, marshy areas, estuaries and inshore coastal waters. The Hovercraft Cruising Club supports the use of hovercraft for cruising in coastal and inland waterways, lakes and lochs. The Hovercraft Club of Great Britain, founded in 1966, regularly organizes inland and coastal hovercraft race events at various venues across the United Kingdom. Similar events are also held in Europe and the US. In August 2010, the Hovercraft Club of Great Britain hosted the World Hovercraft Championships at Towcester Racecourse, followed by the 2016 World Hovercraft Championships at the West Midlands Water Ski Centre in Tamworth. The World Hovercraft Championships are run under the auspices of the World Hovercraft Federation. So far the World Hovercraft Championships had been hosted by France: 1993 in Verneuil, 1997 in Lucon, 2006 at the Lac de Tolerme; Germany: 1987 in Bad Karlshafen, 2004 in Berlin, 2012 and 2018 in Saalburg; Portugal: 1995 in Peso de la Regua; Sweden: 2008 and 2022 at Flottbro Ski Centre in Huddinge; UK 1991 and 2000 at Weston Parc; US: 1989 in Troy (Ohio), 2002 in Terre Haute. The 2020 World Hovercraft Championships had to be postponed to 2022 due to restriction caused by the Covid-19 outbreak. Apart from the craft designed as "racing hovercraft", which are often only suitable for racing, there is another form of small personal hovercraft for leisure use, often referred to as cruising hovercraft, capable of carrying up to four people. Just like their full size counterparts, the ability of these small personal hovercraft to safely cross all types of terrain, (e.g. water, sandbanks, swamps, ice, etc.) and reach places often inaccessible by any other type of craft, makes them suitable for a number of roles, such as survey work and patrol and rescue duties in addition to personal leisure use. Increasingly, these craft are being used as yacht tenders, enabling yacht owners and guests to travel from a waiting yacht to, for example, a secluded beach. In this role, small hovercraft can offer a more entertaining alternative to the usual small boat and can be a rival for the jet-ski. The excitement of a personal hovercraft can now be enjoyed at "experience days", which are popular with families, friends and those in business, who often see them as team building exercises. This level of interest has naturally led to a hovercraft rental sector and numerous manufacturers of small, ready built designs of personal hovercraft to serve the need. Other uses Hoverbarge A real benefit of air cushion vehicles in moving heavy loads over difficult terrain, such as swamps, was overlooked by the excitement of the British Government funding to develop high-speed hovercraft. It was not until the early 1970s that the technology was used for moving a modular marine barge with a dragline on board for use over soft reclaimed land. Mackace (Mackley Air Cushion Equipment), now known as Hovertrans, produced a number of successful Hoverbarges, such as the 250 ton payload "Sea Pearl", which operated in Abu Dhabi, and the twin 160 ton payload "Yukon Princesses", which ferried trucks across the Yukon River to aid the pipeline build. Hoverbarges are still in operation today. In 2006, Hovertrans (formed by the original managers of Mackace) launched a 330-ton payload drilling barge in the swamps of Suriname. The Hoverbarge technology is somewhat different from high-speed hovercraft, which has traditionally been constructed using aircraft technology. The initial concept of the air cushion barge has always been to provide a low-tech amphibious solution for accessing construction sites using typical equipment found in this area, such as diesel engines, ventilating fans, winches and marine equipment. The load to move a 200 ton payload ACV barge at would only be 5 tons. The skirt and air distribution design on high-speed craft again is more complex, as they have to cope with the air cushion being washed out by a wave and wave impact. The slow speed and large mono chamber of the hover barge actually helps reduce the effect of wave action, giving a very smooth ride. The low pull force enabled a Boeing 107 helicopter to pull a hoverbarge across snow, ice and water in 1982. Hovertrains Several attempts have been made to adopt air cushion technology for use in fixed track systems, in order to use the lower frictional forces for delivering high speeds. The most advanced example of this was the Aérotrain, an experimental high speed hovertrain built and operated in France between 1965 and 1977. The project was abandoned in 1977 due to lack of funding, the death of its lead engineer and the adoption of the TGV by the French government as its high-speed ground transport solution. A test track for a tracked hovercraft system was built at Earith near Cambridge, England. It ran southwest from Sutton Gault, sandwiched between the Old Bedford River and the smaller Counter Drain to the west. Careful examination of the site will still reveal traces of the concrete piers used to support the structure. The actual vehicle, RTV31, is preserved at Railworld in Peterborough and can be seen from trains, just south west of Peterborough railway station. The vehicle achieved on 7 February 1973 but the project was cancelled a week later. The project was managed by Tracked Hovercraft Ltd., with Denys Bliss as Director in the early 1970s, then axed by the Aerospace Minister, Michael Heseltine. Records of this project are available from the correspondence and papers of Sir Harry Legge-Bourke, MP at Leeds University Library. Heseltine was accused by Airey Neave and others of misleading the House of Commons when he stated that the government was still considering giving financial support to the Hovertrain, when the decision to pull the plug had already been taken by the Cabinet. After the Cambridge project was abandoned due to financial constraints, parts of the project were picked up by the engineering firm Alfred McAlpine, and abandoned in the mid-1980s. The Tracked Hovercraft project and Professor Laithwaite's Maglev train system were contemporaneous, and there was intense competition between the two prospective British systems for funding and credibility. At the other end of the speed spectrum, the U-Bahn Serfaus has been in continuous operation since 1985. This is an unusual underground air cushion funicular rapid transit system, situated in the Austrian ski resort of Serfaus. Only long, the line reaches a maximum speed of . A similar system also exists in Narita International Airport near Tokyo, Japan. In the late 1960s and early 1970s, the U.S. Department of Transport's Urban Mass Transit Administration (UMTA) funded several hovertrain projects, which were known as Tracked Air Cushion Vehicles or TACVs. They were also known as Aerotrains since one of the builders had a licence from Bertin's Aerotrain company. Three separate projects were funded. Research and development was carried out by Rohr, Inc., Garrett AiResearch and Grumman. UMTA built an extensive test site in Pueblo, Colorado, with different types of tracks for the different technologies used by the prototype contractors. They managed to build prototypes and do a few test runs before the funding was cut. Heavy haulage From the 1960s to 1980s, heavy haulers in the UK used an air-cushion system for their hydraulic modular trailers to carry overweight loads over bridges which were not able to bear the weight of the load and the trailer. The Central Electricity Generating Board had to move transformers from one place to another which weighed from 150 tons to 300 tons for which they did not have appropriate equipment; so they hired heavy haulers like Wynns and Pickfords who had specialized equipment like hydraulic modular trailers manufactured by Nicolas and Cometto, and ballast tractors from Scammell which were strong and powerful enough to carry the load. This made the transportation efficient by avoiding bridge reinforcement, in some cases costing . The transformers were loaded into the girder frame of the hydraulic modular trailer with axle lines in front and behind of the transformer, which made it possible to keep the transformer as low as possible to the ground to negotiate obstacles on the route. Air cushions were mounted under the girder frame's surface and were operated by a compressor vehicle which was a customized Commer 16-ton maxiload provided by CEGB. The vehicle was loaded with 4 air compressors powered by a Rolls-Royce engine producing 235 bhp. While negotiating a bridge the air cushions were inflated and that reduced the stress tremendously on the bridge. Without this technology the government would have had to rebuild the bridges which was not feasible just to carry a small number of loads. Non-transportation The Hoover Constellation was a spherical canister-type vacuum cleaner notable for its lack of wheels. Floating on a cushion of air, it was a domestic hovercraft. They were not especially good as vacuum cleaners as the air escaping from under the cushion blew uncollected dust in all directions, nor as hovercraft as their lack of a skirt meant that they only hovered effectively over a smooth surface. Despite this, original Constellations are sought-after collectibles today. The Flymo is an air-cushion lawn mower that uses a fan on the cutter blade to provide lift. This allows it to be moved in any direction, and provides double-duty as a mulcher. The Marylebone Cricket Club owns a "hover cover" that it uses regularly to cover the pitch at Lord's Cricket Ground. This device is easy and quick to move, and has no pressure points, making damage to the pitch less likely. A power trowel is a hovercraft device used for skimming concrete. Features Advantages Terrain-independence - crossing beachfronts and slopes up to 40 degrees All-season capability - frozen or flowing rivers no object Speed Flexibility, due to low surface friction Disadvantages Engine noise emissions Initial costs Proneness to contrary winds Skirt wear and tear Preservation The Hovercraft Museum at Lee-on-the-Solent, Hampshire, England, houses the world's largest collection of hovercraft designs, including some of the earliest and largest. Much of the collection is housed within the retired SR.N4 hovercraft Princess Anne. She is the last of her kind in the world. There are many hovercraft in the museum but all are non-operational. , Hovercraft continue in use between Ryde on the Isle of Wight and Southsea on the English mainland. The service, operated by Hovertravel, schedules up to three crossings each hour, and provides the fastest way of getting on or off the island. Large passenger-hovercraft are still manufactured on the Isle of Wight. Records World's largest civil hovercraft – The BHC SR.N4 Mk.III, at 56.4 m (185 ft) length and 310 metric tons (305 long tons) weight, can accommodate 418 passengers and 60 cars. World's largest military hovercraft – The Russian Zubr class LCAC at 57.6 metres (188 feet) length and a maximum displacement of 535 tons. This hovercraft can transport three T-80 main battle tanks (MBT), 140 fully equipped troops, or up to 130 tons of cargo. Four have been purchased by the Greek Navy. English Channel crossing – 22 minutes by Princess Anne Mountbatten class hovercraft SR.N4 Mk.III on 14 September 1995 World hovercraft speed record – 137.4 km/h (85.38 mph or 74.19 knots). Bob Windt (USA) at World Hovercraft Championships, Rio Douro River, Peso de Regua, Portugal on 18 September 1995. Hovercraft land speed record – 56.25 mph (90.53 km/h or 48.88 knots). John Alford (USA) at Bonneville Salt Flats, Utah, USA on 21 September 1998. Longest continuous use – The original prototype SR.N6 Mk.I (009) was in service for over 20 years, and logged 22,000 hours of use. It is currently on display at the Hovercraft Museum in Lee-on-the-Solent, Hampshire, England.	Technology	Maritime transport	null
152692	https://en.wikipedia.org/wiki/Tractor	Tractor	A tractor is an engineering vehicle specifically designed to deliver a high tractive effort (or torque) at slow speeds, for the purposes of hauling a trailer or machinery such as that used in agriculture, mining or construction. Most commonly, the term is used to describe a farm vehicle that provides the power and traction to mechanize agricultural tasks, especially (and originally) tillage, and now many more. Agricultural implements may be towed behind or mounted on the tractor, and the tractor may also provide a source of power if the implement is mechanised. Etymology The word tractor was taken from Latin, being the agent noun of trahere "to pull". The first recorded use of the word meaning "an engine or vehicle for pulling wagons or plows" occurred in 1896, from the earlier term "traction motor" (1859). National variations In the UK, Ireland, Australia, India, Spain, Argentina, Slovenia, Serbia, Croatia, the Netherlands, and Germany, the word "tractor" usually means "farm tractor", and the use of the word "tractor" to mean other types of vehicles is familiar to the vehicle trade, but unfamiliar to much of the general public. In Canada and the US, the word may also refer to the road tractor portion of a tractor trailer truck, but also usually refers to the piece of farm equipment. History Traction engines The first powered farm implements in the early 19th century were portable engines – steam engines on wheels that could be used to drive mechanical farm machinery by way of a flexible belt. Richard Trevithick designed the first 'semi-portable' stationary steam engine for agricultural use, known as a "barn engine" in 1812, and it was used to drive a corn threshing machine. The truly portable engine was invented in 1839 by William Tuxford of Boston, Lincolnshire who started manufacture of an engine built around a locomotive-style boiler with horizontal smoke tubes. A large flywheel was mounted on the crankshaft, and a stout leather belt was used to transfer the drive to the equipment being driven. In the 1850s, John Fowler used a Clayton & Shuttleworth portable engine to drive apparatus in the first public demonstrations of the application of cable haulage to cultivation. In parallel with the early portable engine development, many engineers attempted to make them self-propelled – the fore-runners of the traction engine. In most cases this was achieved by fitting a sprocket on the end of the crankshaft, and running a chain from this to a larger sprocket on the rear axle. These experiments met with mixed success. The first proper traction engine, in the form recognisable today, was developed in 1859 when British engineer Thomas Aveling modified a Clayton & Shuttleworth portable engine, which had to be hauled from job to job by horses, into a self-propelled one. The alteration was made by fitting a long driving chain between the crankshaft and the rear axle. The first half of the 1860s was a period of great experimentation but by the end of the decade the standard form of the traction engine had evolved and changed little over the next sixty years. It was widely adopted for agricultural use. The first tractors were steam-powered plowing engines. They were used in pairs, placed on either side of a field to haul a plow back and forth between them using a wire cable. In Britain Mann's and Garrett developed steam tractors for direct ploughing, but the heavy, wet soil of England meant that these designs were less economical than a team of horses. In the United States, where soil conditions permitted, steam tractors were used to direct-haul plows. Steam-powered agricultural engines remained in use well into the 20th century until reliable internal combustion engines had been developed. Fuel The first gasoline powered tractors were built in Illinois, by John Charter combining single cylinder Otto engines with a Rumley Steam engine chassis, in 1889. In 1892, John Froelich built a gasoline-powered tractor in Clayton County, Iowa, US. A Van Duzen single-cylinder gasoline engine was mounted on a Robinson engine chassis, which could be controlled and propelled by Froelich's gear box. After receiving a patent, Froelich started up the Waterloo Gasoline Engine Company and invested all of his assets. The venture was very unsuccessful, and by 1895 all was lost and he went out of business. Richard Hornsby & Sons are credited with producing and selling the first oil-engined tractor in Britain, invented by Herbert Akroyd Stuart. The Hornsby-Akroyd Patent Safety Oil Traction Engine was made in 1896 with a engine. In 1897, it was bought by Mr. Locke-King, the first recorded British tractor sale. That year, it won a Silver Medal from the Royal Agricultural Society of England. It later returned to the factory for a caterpillar track fitting. The first commercially successful light-weight petrol-powered general purpose tractor was built by Dan Albone, a British inventor in 1901. He filed for a patent on 15 February 1902 for his tractor design and then formed Ivel Agricultural Motors Limited. The other directors were Selwyn Edge, Charles Jarrott, John Hewitt and Lord Willoughby. He called his machine the Ivel Agricultural Motor; the word "tractor" came into common use after Hart-Parr created it. The Ivel Agricultural Motor was light, powerful and compact. It had one front wheel, with a solid rubber tyre, and two large rear wheels like a modern tractor. The engine used water cooling, utilizing the thermo-syphon effect. It had one forward and one reverse gear. A pulley wheel on the left hand side allowed it to be used as a stationary engine, driving a wide range of agricultural machinery. The 1903 sale price was £300. His tractor won a medal at the Royal Agricultural Show, in 1903 and 1904. About 500 were built, and many were exported all over the world. The original engine was made by Payne & Co. of Coventry. After 1906, French Aster engines were used. The first successful American tractor was built by Charles W. Hart and Charles H. Parr. They developed a two-cylinder gasoline engine and set up their business in Charles City, Iowa. In 1903, the firm built 15 tractors. Their #3 is the oldest surviving internal combustion engine tractor in the United States, and is on display at the Smithsonian National Museum of American History in Washington, D.C. The two-cylinder engine has a unique hit-and-miss firing cycle that produced at the belt and at the drawbar. In 1908, the Saunderson Tractor and Implement Co. of Bedford introduced a four-wheel design, and became the largest tractor manufacturer in Britain at the time. While the earlier, heavier tractors were initially very successful, it became increasingly apparent at this time that the weight of a large supporting frame was less efficient than lighter designs. Henry Ford introduced a light-weight, mass-produced design which largely displaced the heavier designs. Some companies halfheartedly followed suit with mediocre designs, as if to disprove the concept, but they were largely unsuccessful in that endeavor. While unpopular at first, these gasoline-powered machines began to catch on in the 1910s, when they became smaller and more affordable. Henry Ford introduced the Fordson, a wildly popular mass-produced tractor, in 1917. They were built in the U.S., Ireland, England and Russia, and by 1923, Fordson had 77% of the U.S. market. The Fordson dispensed with a frame, using the strength of the engine block to hold the machine together. By the 1920s, tractors with gasoline-powered internal combustion engines had become the norm. The first three-point hitches were experimented with in 1917. After Harry Ferguson applied for a British patent for his three-point hitch in 1926, they became popular. A three-point attachment of the implement to the tractor is the simplest and the only statically determinate way of joining two bodies in engineering. The Ferguson-Brown Company produced the Model A Ferguson-Brown tractor with a Ferguson-designed hydraulic hitch. In 1938 Ferguson entered into a collaboration with Henry Ford to produce the Ford-Ferguson 9N tractor. The three-point hitch soon became the favorite hitch attachment system among farmers around the world. This tractor model also included a rear Power Take Off (PTO) shaft that could be used to power three point hitch mounted implements such as sickle-bar mowers. Electric In 1969, General Electric introduced the Elec-Trak, the first commercial, electric tractor (electric-powered garden tractor). The Elec-Trak was manufactured by General Electric until 1975. Electric tractors are manufactured by a German company, Fendt, and by US companies, Solectrac and Monarch Tractor. John Deere's protoype electric tractor is a plug-in, powered by an electrical cable. Kubota is prototyping an autonomous electric tractor. Design, power and transmission Configuration Tractors can be generally classified by number of axles or wheels, with main categories of two-wheel tractors (single-axle tractors) and four-wheel tractors (two-axle tractors); more axles are possible but uncommon. Among four-wheel tractors (two-axle tractors), most are two-wheel drive (usually at the rear); but many are two-wheel drive with front wheel assist, four-wheel drive (often with articulated steering), or track crawler (with steel or rubber tracks). The classic farm tractor is a simple open vehicle, with two very large driving wheels on an axle below a single seat (the seat and steering wheel consequently are in the center), and the engine in front of the driver, with two steerable wheels below the engine compartment. This basic design has remained unchanged for a number of years after being pioneered by Wallis, but enclosed cabs are fitted on almost all modern models, for operator safety and comfort. In some localities with heavy or wet soils, notably in the Central Valley of California, the "Caterpillar" or "crawler" type of tracked tractor became popular due to superior traction and flotation. These were usually maneuvered through the use of turning brake pedals and separate track clutches operated by levers rather than a steering wheel. Four-wheel drive tractors began to appear in the 1960s. Some four-wheel drive tractors have the standard "two large, two small" configuration typical of smaller tractors, while some have four large, powered wheels. The larger tractors are typically an articulated, center-hinged design steered by hydraulic cylinders that move the forward power unit while the trailing unit is not steered separately. In the early 21st century, articulated or non-articulated, steerable multitrack tractors have largely supplanted the Caterpillar type for farm use. Larger types of modern farm tractors include articulated four-wheel or eight-wheel drive units with one or two power units which are hinged in the middle and steered by hydraulic clutches or pumps. A relatively recent development is the replacement of wheels or steel crawler-type tracks with flexible, steel-reinforced rubber tracks, usually powered by hydrostatic or completely hydraulic driving mechanisms. The configuration of these tractors bears little resemblance to the classic farm tractor design. Engine and fuels The predecessors of modern tractors, traction engines, used steam engines for power. Gasoline and kerosene Since the turn of the 20th century, internal combustion engines have been the power source of choice. Between 1900 and 1960, gasoline was the predominant fuel, with kerosene (the Rumely Oil Pull was the most notable of this kind)being a common alternative. Generally, one engine could burn any of those, although cold starting was easiest on gasoline. Often, a small auxiliary fuel tank was available to hold gasoline for cold starting and warm-up, while the main fuel tank held whatever fuel was most convenient or least expensive for the particular farmer. In the United Kingdom, a gasoline-kerosene engine is known as a petrol-paraffin engine. Diesel Dieselisation gained momentum starting in the 1960s, and modern farm tractors usually employ diesel engines, which range in power output from 18 to 575 horsepower (15 to 480 kW). Size and output are dependent on application, with smaller tractors used for lawn mowing, landscaping, orchard work, and truck farming, and larger tractors for vast fields of wheat, corn, soy, and other bulk crops. Liquefied petroleum gas Liquefied petroleum gas (LPG) or propane also have been used as tractor fuels, but require special pressurized fuel tanks and filling equipment and produced less power, so are less prevalent in most markets. Most are confined for inside work due to their clean burning. Wood During the second world war, Petrolium based fuel was scarce in many European nations. So they resorted to using wood gasifires on every vehicle, including tractors. Biodiesel In some countries such as Germany, biodiesel is often used. Some other biofuels such as straight vegetable oil are also being used by some farmers. Electric powered Prototype battery powered electric tractors are being developed by a German company, Fendt, and by two US companies, Solectrac and Monarch Tractor. John Deere's protoype electric tractor is a plug-in, powered by an electrical cable. Kubota is prototyping an autonomous electric tractor. Transmission Most older farm tractors use a manual transmission with several gear ratios, typically three to six, sometimes multiplied into two or three ranges. This arrangement provides a set of discrete ratios that, combined with the varying of the throttle, allow final-drive speeds from less than one up to about 25 miles per hour (40 km/h), with the lower speeds used for working the land and the highest speed used on the road. Slow, controllable speeds are necessary for most of the operations performed with a tractor. They help give the farmer a larger degree of control in certain situations, such as field work. When travelling on public roads, the slow operating speeds can cause problems, such as long queues or tailbacks, which can delay or annoy motorists in cars and trucks. These motorists are responsible for being duly careful around farm tractors and sharing the road with them, but many shirk this responsibility, so various ways to minimize the interaction or minimize the speed differential are employed where feasible. Some countries (for example the Netherlands) employ a road sign on some roads that means "no farm tractors". Some modern tractors, such as the JCB Fastrac, are now capable of much higher road speeds of around 50 mph (80 km/h). Older tractors usually have unsynchronized transmission designs, which often require the operator to engage the clutch to shift between gears. This mode of use is inherently unsuited to some of the work tractors do, and has been circumvented in various ways over the years. For existing unsynchronized tractors, the methods of circumvention are double clutching or power-shifting, both of which require the operator to rely on skill to speed-match the gears while shifting, and are undesirable from a risk-mitigation standpoint because of what can go wrong if the operator makes a mistake – transmission damage is possible, and loss of vehicle control can occur if the tractor is towing a heavy load either uphill or downhill – something that tractors often do. Therefore, operator's manuals for most of these tractors state one must always stop the tractor before shifting. In newer designs, unsynchronized transmission designs were replaced with synchronization or with continuously variable transmissions (CVTs). Either a synchronized manual transmission with enough available gear ratios (often achieved with dual ranges, high and low) or a CVT allow the engine speed to be matched to the desired final-drive speed, while keeping engine speed within the appropriate speed (as measured in rotations per minute or rpm) range for power generation (the working range) (whereas throttling back to achieve the desired final-drive speed is a trade-off that leaves the working range). The problems, solutions, and developments described here also describe the history of transmission evolution in semi-trailer trucks. The biggest difference is fleet turnover; whereas most of the old road tractors have long since been scrapped, many of the old farm tractors are still in use. Therefore, old transmission design and operation is primarily just of historical interest in trucking, whereas in farming it still often affects daily life. Hitches and power applications The power produced by the engine must be transmitted to the implement or equipment to do the actual work intended for the equipment. This may be accomplished via a drawbar or hitch system if the implement is to be towed or otherwise pulled through the tractive power of the engine, or via a pulley or power takeoff system if the implement is stationary, or a combination of the two. Drawbars Plows and other tillage equipment are most commonly connected to the tractor via a drawbar. The classic drawbar is simply a steel bar attached to the tractor (or in some cases, as in the early Fordsons, cast as part of the rear transmission housing) to which the hitch of the implement was attached with a pin or by a loop and clevis. The implement could be readily attached and removed, allowing the tractor to be used for other purposes on a daily basis. If the tractor was equipped with a swinging drawbar, then it could be set at the center or offset from center to allow the tractor to run outside the path of the implement. The drawbar system necessitated the implement having its own running gear (usually wheels) and in the case of a plow, chisel cultivator or harrow, some sort of lift mechanism to raise it out of the ground at turns or for transport. Drawbars necessarily posed a rollover risk depending on how the tractive torque was applied. The Fordson tractor was prone to roll backward due to an excessively short wheelbase. The linkage between the implement and the tractor usually had some slack which could lead to jerky starts and greater wear and tear on the tractor and the equipment. Drawbars were appropriate to the dawn of mechanization, because they were very simple in concept and because as the tractor replaced the horse, existing horse-drawn implements usually already had running gear. As the history of mechanization progressed, the advantages of other hitching systems became apparent, leading to new developments (see below). Depending on the function for which a tractor is used, though, the drawbar is still one of the usual means of attaching an implement to a tractor (see photo at left). Fixed mounts Some tractor manufacturers produced matching equipment that could be directly mounted on the tractor. Examples included front-end loaders, belly mowers, row crop cultivators, corn pickers and corn planters. In most cases, these fixed mounts were proprietary and unique to each make of tractor, so an implement produced by John Deere, for example, could not be attached to a Minneapolis Moline tractor. Another disadvantage was mounting usually required some time and labor, resulting in the implement being semi-permanently attached with bolts or other mounting hardware. Usually, it was impractical to remove the implement and reinstall it on a day-to-day basis. As a result, the tractor was unavailable for other uses and dedicated to a single use for an appreciable period of time. An implement was generally mounted at the beginning of its season of use (such as tillage, planting or harvesting) and removed when the season ended. Three-point and quick The drawbar system was virtually the exclusive method of attaching implements (other than direct attachment to the tractor) before Harry Ferguson developed the three-point hitch. Equipment attached to the three-point hitch can be raised or lowered hydraulically with a control lever. The equipment attached to the three-point hitch is usually completely supported by the tractor. Another way to attach an implement is via a quick hitch, which is attached to the three-point hitch. This enables a single person to attach an implement quicker and put the person in less danger when attaching the implement. The three-point hitch revolutionized farm tractors and their implements. While the Ferguson System was still under patent, other manufacturers developed new hitching systems to try to fend off some of Ferguson's competitive advantage. For example, International Harvester's Farmall tractors gained a two-point "Fast Hitch", and John Deere had a power lift that was somewhat similar to the more flexible Ferguson invention. Once the patent protection expired on the three-point hitch, it became an industry standard. Almost every tractor today features Ferguson's three-point linkage or a derivative of it. This hitch allows for easy attachment and detachment of implements while allowing the implement to function as a part of the tractor, almost as if it were attached by a fixed mount. Previously, when the implement hit an obstacle, the towing link broke or the tractor flipped over. Ferguson's idea was to combine a connection via two lower and one upper lift arms that were connected to a hydraulic lifting ram. The ram was, in turn, connected to the upper of the three links so the increased drag (as when a plough hits a rock) caused the hydraulics to lift the implement until the obstacle was passed. Recently, Bobcat's patent on its front loader connection (inspired by these earlier systems) has expired, and compact tractors are now being outfitted with quick-connect attachments for their front-end loaders. Power take-off systems and hydraulics In addition to towing an implement or supplying tractive power through the wheels, most tractors have a means to transfer power to another machine such as a baler, swather, or mower. Unless it functions solely by pulling it through or over the ground, a towed implement needs its own power source (such as a baler or combine with a separate engine) or else a means of transmitting power from the tractor to the mechanical operations of the equipment. Early tractors used belts or cables wrapped around the flywheel or a separate belt pulley to power stationary equipment, such as a threshing machine, buzz saw, silage blower, or stationary baler. In most cases, it was impractical for the tractor and equipment to move with a flexible belt or cable between them, so this system required the tractor to remain in one location, with the work brought to the equipment, or the tractor to be relocated at each turn and the power set-up reapplied (as in cable-drawn plowing systems used in early steam tractor operations). Modern tractors use a power take-off (PTO) shaft to provide rotary power to machinery that may be stationary or pulled. The PTO shaft generally is at the rear of the tractor, and can be connected to an implement that is either towed by a drawbar or a three-point hitch. This eliminates the need for a separate, implement-mounted power source, which is almost never seen in modern farm equipment. It is also optional to get a front PTO as well when buying a new tractor. Virtually all modern tractors can also provide external hydraulic fluid and electrical power to the equipment they are towing, either by hoses or wires. Operation Modern tractors have many electrical switches and levers in the cab for controlling the multitude of different functions available on the tractor. Pedals Some modern farm tractors retain a traditional manual transmission; increasingly they have hydraulically driven powershift transmissions and CVT, which vastly simplify operation. Those with powershift transmissions have identical pedal arrangements on the floor for the operator to actuate, replacing a clutch pedal on the far left with an inching pedal that cuts off hydraulic flow to the clutches. Twinned brake pedals – one each for left and right side wheels – are placed together on the right side. Some have a pedal for a foot throttle on the far right. Unlike automobiles, throttle speed can also be controlled by a hand-operated lever ("hand throttle"), which may be set to a fixed position. This helps provide a constant speed in field work. It also helps provide continuous power for stationary tractors that are operating an implement by PTO shaft or axle driven belt. The foot throttle gives the operator more automobile-like control over the speed of a mobile tractor in any operation. Some modern tractors also have (or offer as optional equipment) a button on the gear stick for controlling the clutch, in addition to the standard pedal, allowing for gear changes and the tractor to be brought to a stop without using the foot pedal to engage the clutch. Others have a button for temporarily increasing throttle speed to improve hydraulic flow to implements, such as a front end loader bucket. Independent left and right brake pedals are provided to allow improved steering (by engaging the side one wishes to turn to, slowing or stopping its wheel) and improved traction in soft and slippery conditions (by transferring rotation to the wheel with better grip). Some users prefer to lock both pedals together, or utilize a partial lock that allows the left pedal to be depressed independently but engages both when the right is applied. This may be in the form of a swinging or sliding bolt that may be readily engaged or disengaged in the field without tools. Foot pedal throttle control is mostly a returning feature of newer tractors. In the UK, foot pedal use to control engine speed while travelling on the road is mandatory. Some tractors, especially those designed for row-crop work, have a 'de-accelerator' pedal, which operates in the reverse fashion of an automobile throttle, slowing the engine when applied. This allows control over the speed of a tractor with its throttle set high for work, as when repeatedly slowing to make U-turns at the end of crop rows in fields. A front-facing foot button is traditionally included just ahead of the driver's seat (designed to be pressed by the operator's heel) to engage the rear differential lock (diff-lock), which prevents wheel slip. The differential normally allows driving wheels to operate at their own speeds, as required, for example, by the different radius each takes in a turn. This allows the outside wheel to travel faster than the inside wheel, thereby traveling further during a turn. In low-traction conditions on a soft surface, the same mechanism can allow one wheel to slip, wasting its torque and further reducing traction. The differential lock overrides this, forcing both wheels to turn at the same speed, reducing wheel slip and improving traction. Care must be taken to unlock the differential before turning, usually by hitting the pedal a second time, since the tractor with good traction cannot perform a turn with the diff-lock engaged. In many modern tractors, this pedal is replaced with an electrical switch. Levers and switches Many functions once controlled with levers have been replaced with some model of electrical switch with the rise of indirect computer controlling of functions in modern tractors. Until the late of the 1950s, tractors had a single register of gears, hence one gear stick, often with three to five forward gears and one reverse. Then, group gears were introduced, and another gear stick was added. Later, control of the forward-reverse direction was moved to a special stick attached at the side of the steering wheel, which allowed forward or reverse travel in any gear. Now, with CVTs or other gear types, fewer sticks control the transmission, and some are replaced with electrical switches or are totally computer-controlled. The three-point hitch was controlled with a lever for adjusting the position, or as with the earliest ones, just the function for raising or lowering the hitch. With modern electrical systems, it is often replaced with a potentiometer for the lower bound position and another one for the upper bound, and a switch allowing automatic adjustment of the hitch between these settings. The external hydraulics also originally had levers, but now are often replaced with some form of electrical switch; the same is true for the power take-off shaft. Safety Agriculture in the United States is one of the most hazardous industries, only surpassed by mining and construction. No other farm machine is so identified with the hazards of production agriculture as the tractor. Tractor-related injuries account for approximately 32% of the fatalities and 6% of the nonfatal injuries in agriculture. Over 50% is attributed to tractor overturns. The roll-over protection structure (ROPS) and seat belt, when worn, are the most important safety devices to protect operators from death during tractor overturns. Modern tractors have a ROPS to prevent an operator from being crushed when overturning. This is especially important in open-air tractors, where the ROPS is a steel beam that extends above the operator's seat. For tractors with operator cabs, the ROPS is part of the frame of the cab. A ROPS with enclosed cab further reduces the likelihood of serious injury because the operator is protected by the sides and windows of the cab. These structures were first required by legislation in Sweden in 1959. Before they were required, some farmers died when their tractors rolled on top of them. Row-crop tractors, before ROPS, were particularly dangerous because of their 'tricycle' design with the two front wheels spaced close together and angled inward toward the ground. Some farmers were killed by rollovers while operating tractors along steep slopes. Others have been killed while attempting to tow or pull an excessive load from above axle height, or when cold weather caused the tires to freeze to the ground, in both cases causing the tractor to pivot around the rear axle. ROPS were first required in the United States in 1986, non-retroactively. ROPS adoption by farmers is thus incomplete. To treat this problem, CROPS (cost-effective roll-over protection structures) have been developed to encourage farmers to retrofit older tractors. For the ROPS to work as designed, the operator must stay within its protective frame and wear the seat belt. In addition to ROPS, U.S. manufacturers add instructional seats on tractors with enclosed cabs. The tractors have a ROPS with seatbelts for both the operator and passenger. This instructional seat is intended to be used for training new tractor operators, but can also be used to diagnose machine problems. The misuse of an instructional seat increases the likelihood of injury, especially when children are transported. The International Organization for Standardization's ISO standard 23205:2014 specifies the minimum design and performance requirements for an instructional seat and states that the instructional seat is neither intended for, nor is it designed for use by children. Despite this, upwards of 40% of farm families give their children rides on tractors, often using these instructional seats. Applications and variations Farm The most common use of the term "tractor" is for the vehicles used on farms. The farm tractor is used for pulling or pushing agricultural machinery or trailers, for plowing, tilling, disking, harrowing, planting, and similar tasks. A variety of specialty farm tractors have been developed for particular uses. These include "row crop" tractors with adjustable tread width to allow the tractor to pass down rows of cereals, maize, tomatoes or other crops without crushing the plants, "wheatland" or "standard" tractors with fixed wheels and a lower center of gravity for plowing and other heavy field work for broadcast crops, and "high crop" tractors with adjustable tread and increased ground clearance, often used in the cultivation of cotton and other high-growing row crop plant operations, and "utility tractors", typically smaller tractors with a low center of gravity and short turning radius, used for general purposes around the farmstead. Many utility tractors are used for nonfarm grading, landscape maintenance and excavation purposes, particularly with loaders, backhoes, pallet forks and similar devices. Small garden or lawn tractors designed for suburban and semirural gardening and landscape maintenance are produced in a variety of configurations, and also find numerous uses on a farmstead. Some farm-type tractors are found elsewhere than on farms: with large universities' gardening departments, in public parks, or for highway workman use with blowtorch cylinders strapped to the sides and a pneumatic drill air compressor permanently fastened over the power take-off. These are often fitted with grass (turf) tyres which are less damaging to soft surfaces than agricultural tires. Precision Space technology has been incorporated into agriculture in the form of GPS devices, and robust on-board computers installed as optional features on farm tractors. These technologies are used in modern, precision farming techniques. The spin-offs from the space race have actually facilitated automation in plowing and the use of autosteer systems (drone on tractors that are manned but only steered at the end of a row), the idea being to neither overlap and use more fuel nor leave streaks when performing jobs such as cultivating. Several tractor companies have also been working on producing a driverless tractor. Engineering The durability and engine power of tractors made them very suitable for engineering tasks. Tractors can be fitted with engineering tools such as dozer blades, buckets, hoes, rippers, etc. The most common attachments for the front of a tractor are dozer blades or buckets. When attached to engineering tools, the tractor is called an engineering vehicle. A bulldozer is a track-type tractor with a blade attached in the front and a rope-winch behind. Bulldozers are very powerful tractors and have excellent ground-hold, as their main tasks are to push or drag. Bulldozers have been further modified over time to evolve into new machines which are capable of working in ways that the original bulldozer can not. One example is that loader tractors were created by removing the blade and substituting a large volume bucket and hydraulic arms which can raise and lower the bucket, thus making it useful for scooping up earth, rock and similar loose material to load it into trucks. A front-loader or loader is a tractor with an engineering tool which consists of two hydraulic powered arms on either side of the front engine compartment and a tilting implement. This is usually a wide-open box called a bucket, but other common attachments are a pallet fork and a bale grappler. Other modifications to the original bulldozer include making the machine smaller to let it operate in small work areas where movement is limited. Also, tiny wheeled loaders, officially called skid-steer loaders, but nicknamed "Bobcat" after the original manufacturer, are particularly suited for small excavation projects in confined areas. Backhoe The most common variation of the classic farm tractor is the backhoe, also called a backhoe-loader. As the name implies, it has a loader assembly on the front and a backhoe on the back. Backhoes attach to a three-point hitch on farm or industrial tractors. Industrial tractors are often heavier in construction, particularly with regards to the use of a steel grill for protection from rocks and the use of construction tires. When the backhoe is permanently attached, the machine usually has a seat that can swivel to the rear to face the hoe controls. Removable backhoe attachments almost always have a separate seat on the attachment. Backhoe-loaders are very common and can be used for a wide variety of tasks: construction, small demolitions, light transportation of building materials, powering building equipment, digging holes, loading trucks, breaking asphalt and paving roads. Some buckets have retractable bottoms, enabling them to empty their loads more quickly and efficiently. Buckets with retractable bottoms are also often used for grading and scratching off sand. The front assembly may be a removable attachment or permanently mounted. Often the bucket can be replaced with other devices or tools. Their relatively small frames and precise controls make backhoe-loaders very useful and common in urban engineering projects, such as construction and repairs in areas too small for larger equipment. Their versatility and compact size make them one of the most popular urban construction vehicles. In the UK and Ireland, the word "JCB" is used colloquially as a genericized trademark for any such type of engineering vehicle. The term JCB now appears in the Oxford English Dictionary, although it is still legally a trademark of J. C. Bamford Ltd. The term "digger" is also commonly used. Compact utility A compact utility tractor (CUT) is a smaller version of an agricultural tractor, but designed primarily for landscaping and estate management tasks, rather than for planting and harvesting on a commercial scale. Typical CUTs range from with available power take-off (PTO) power ranging from . CUTs are often equipped with both a mid-mounted and a standard rear PTO, especially those below . The mid-mount PTO shaft typically rotates at/near 2000 rpm and is typically used to power mid-mount finish mowers, front-mounted snow blowers or front-mounted rotary brooms. The rear PTO is standardized at 540 rpm for the North American markets, but in some parts of the world, a dual 540/1000 rpm PTO is standard, and implements are available for either standard in those markets. One of the most common attachments for a CUT is the front-end loader or FEL. Like the larger agricultural tractors, a CUT will have an adjustable, hydraulically controlled three-point hitch. Typically, a CUT will have four-wheel drive, or more correctly four-wheel assist. Modern CUTs often feature hydrostatic transmissions, but many variants of gear-drive transmissions are also offered from low priced, simple gear transmissions to synchronized transmissions to advanced glide-shift transmissions. All modern CUTs feature government-mandated roll over protection structures just like agricultural tractors. The most well-known brands in North America include Kubota, John Deere Tractor, New Holland Ag, Case-Farmall and Massey Ferguson. Although less common, compact backhoes are often attached to compact utility tractors. Compact utility tractors require special, smaller implements than full-sized agricultural tractors. Very common implements include the box blade, the grader blade, the landscape rake, the post hole digger (or post hole auger), the rotary cutter (slasher or a brush hog), a mid- or rear-mount finish mower, a broadcast seeder, a subsoiler and the rototiller (rotary tiller). In northern climates, a rear-mounted snow blower is very common; some smaller CUT models are available with front-mounted snow blowers powered by mid-PTO shafts. Implement brands outnumber tractor brands, so CUT owners have a wide selection of implements. For small-scale farming or large-scale gardening, some planting and harvesting implements are sized for CUTs. One- and two-row planting units are commonly available, as are cultivators, sprayers and different types of seeders (slit, rotary and drop). One of the first CUTs offered for small farms of three to 30 acres and for small jobs on larger farms was a three-wheeled unit, with the rear wheel being the drive wheel, offered by Sears & Roebuck in 1954 and priced at $598 for the basic model. An even smaller variant of the compact utility tractor is the subcompact utility tractor. Although these tractors are often barely larger than a riding lawn mower, these tractors have all the same features of a compact tractor, such as a three-point hitch, power steering, four-wheel-drive, and front-end loader. These tractors are generally marketed towards homeowners who intend to mostly use them for lawn mowing, with the occasional light landscaping task. Standard The earliest tractors were called "standard" tractors, and were intended almost solely for plowing and harrowing before planting, which were difficult tasks for humans and draft animals. They were characterized by a low, rearward seating position, fixed-width tread, and low ground clearance. These early tractors were cumbersome, and ill-suited to enter a field of planted row crops for weed control. The "standard" tractor definition is no longer in current use. However, tractors with fixed wheel spacing and a low center of gravity are well-suited as loaders, forklifts and backhoes, so that the configuration continues in use without the "standard" nomenclature. Row-crop A general-purpose or row-crop tractor is tailored specifically to the growing of crops grown in rows, and most especially to cultivating these crops. These tractors are universal machines, capable of both primary tillage and cultivation of a crop. The row-crop tractor category evolved rather than appearing overnight, but the International Harvester (IH) Farmall is often considered the "first" tractor of the category. Some earlier tractors of the 1910s and 1920s approached the form factor from the heavier side, as did motorized cultivators from the lighter side, but the Farmall brought all of the salient features together into one package, with a capable distribution network to ensure its commercial success. In the new form factor that the Farmall popularized, the cultivator was mounted in the front so it was easily visible. Additionally, the tractor had a narrow front end; the front tires were spaced very closely and angled in toward the bottom. The back wheels straddled two rows with their spacing adjustable depending on row spacing, and the unit could cultivate four rows at once. Where wide front wheels were used, they often could be adjusted as well. Tractors with non-adjustable spacing were called "standard" or "wheatland", and were chiefly meant for pulling plows or other towed implements, typically with a lower overall tractor height than row-crop models. From 1924 until 1963, Farmalls were the largest selling row-crop tractors. To compete, John Deere designed the Model C, which had a wide front and could cultivate three rows at once. Only 112 prototypes were made, as Deere realized it would lose sales to Farmall if its model did less. In 1928, Deere released the Model C anyway, only as the Model GP (General Purpose) to avoid confusion with the Model D when ordered over the then unclear telephone. Oliver refined its "Row Crop" model early in 1930. Until 1935, the 18–27 was Oliver–Hart-Parr's only row-crop tractor. Many Oliver row-crop models are referred to as "Oliver Row Crop 77", "Oliver Row Crop 88", etc. Many early row-crop tractors had a tricycle design with two closely spaced front tires, and some even had a single front tire. This made it dangerous to operate on the side of a steep hill; as a result, many farmers died from tractor rollovers. Also, early row-crop tractors had no rollover protection system (ROPS), meaning if the tractor flipped back, the operator could be crushed. Sweden was the first country which passed legislation requiring ROPS, in 1959. Over 50% of tractor related injuries and deaths are attributed to tractor rollover. Canadian agricultural equipment manufacturer Versatile makes row-crop tractors that are ; powered by an 8.9 liter Cummins Diesel engine. Case IH and New Holland of CNH Industrial both produce high horsepower front-wheel-assist row crop tractors with available rear tracks. Case IH also has a four-wheel drive track system called Rowtrac. John Deere has an extensive line of available row crop tractors ranging from . Modern row crop tractors have rollover protection systems in the form of a reinforced cab or a roll bar. Garden Garden tractors, sometimes called lawn tractors, are small, light tractors designed for use in domestic gardens, lawns, and small estates. Lawn tractors are designed for cutting grass and snow removal, while garden tractors are for small property cultivation. In the U.S., the term riding lawn mower today often is used to refer to mid- or rear-engined machines. Front-engined tractor layout machines designed primarily for cutting grass and light towing are called lawn tractors; heavier-duty tractors of similar size are garden tractors. Garden tractors are capable of mounting a wider array of attachments than lawn tractors. Unlike lawn tractors and rear-engined riding mowers, garden tractors are powered by horizontal-crankshaft engines with a belt-drive to transaxle-type transmissions (usually of four or five speeds, although some may also have two-speed reduction gearboxes, drive-shafts, or hydrostatic or hydraulic drives). Garden tractors from Wheel Horse, Cub Cadet, Economy (Power King), John Deere, Massey Ferguson and Case Ingersoll are built in this manner. The engines are generally one- or two-cylinder petrol (gasoline) engines, although diesel engine models are also available, especially in Europe. Typically, diesel-powered garden tractors are larger and heavier-duty than gasoline-powered units and compare more similarly to compact utility tractors. Visually, the distinction between a garden tractor and a lawn tractor is often hard to make – generally, garden tractors are more sturdily built, with stronger frames, 12-inch or larger wheels mounted with multiple lugs (most lawn tractors have a single bolt or clip on the hub), heavier transaxles, and ability to accommodate a wide range of front, belly, and rear mounted attachments. Two-wheel Although most people think primarily of four-wheel vehicles when they think of tractors, a tractor may have one or more axles. The key benefit is the power itself, which only takes one axle to provide. Single-axle tractors, more often called two-wheel tractors or walk-behind tractors, have had many users since the introduction of the internal combustion engine tractors. They tend to be small and affordable, this was especially true before the 1960s when a walk-behind tractor could often be more affordable than a two-axle tractor of comparable power. Today's compact utility tractors and advanced garden tractors may negate most of that market advantage, but two-wheel tractors still have a following, especially among those who already own one. Countries where two-wheel tractors are especially prevalent today include Thailand, China, Bangladesh, India, and other Southeast Asia countries. Most two-wheel tractors today are specialty tractors made for one purpose, such as snow blowers, push tillers, and self propelled push mowers. Orchard Tractors tailored to use in fruit orchards typically have features suited to passing under tree branches with impunity. These include a lower overall profile; reduced tree-branch-snagging risk (via underslung exhaust pipes rather than smoke-stack-style exhaust, and large sheetmetal cowlings and fairings that allow branches to deflect and slide off rather than catch); spark arrestors on the exhaust tips; and often wire cages to protect the operator from snags. Automobile conversions and other homemade versions The ingenuity of farm mechanics, coupled in some cases with OEM or aftermarket assistance, has often resulted in the conversion of automobiles for use as farm tractors. In the United States, this trend was especially strong from the 1910s through 1950s. It began early in the development of vehicles powered by internal combustion engines, with blacksmiths and amateur mechanics tinkering in their shops. Especially during the interwar period, dozens of manufacturers (Montgomery Ward among them) marketed aftermarket kits for converting Ford Model Ts for use as tractors. (These were sometimes called 'Hoover wagons' during the Great Depression, although this term was usually reserved for automobiles converted to horse-drawn buggy use when gasoline was unavailable or unaffordable. During the same period, another common name was "Doodlebug", after the popular kit by the same name.) Ford even considered producing an "official" optional kit. Many Model A Fords also were converted for this purpose. In later years, some farm mechanics have been known to convert more modern trucks or cars for use as tractors, more often as curiosities or for recreational purposes (rather than out of the earlier motives of pure necessity or frugality). During World War II, a shortage of tractors in Sweden led to the development of the so-called "EPA" tractor (EPA was a chain of discount stores and it was often used to signify something lacking in quality). An EPA tractor was simply an automobile, truck or lorry, with the passenger space cut off behind the front seats, equipped with two gearboxes in a row. When done to an older car with a ladder frame, the result was similar to a tractor and could be used as one. After the war it remained popular as a way for young people without a driver's license to own something similar to a car. Since it was legally seen as a tractor, it could be driven from 16 years of age and only required a tractor license. Eventually, the legal loophole was closed and no new EPA tractors were allowed to be made, but the remaining ones were still legal, which led to inflated prices and many protests from people who preferred EPA tractors to ordinary cars. The Swedish government eventually replaced them with the so called "A-tractor" which now had its speed limited to 30 km/h and allowed people aged 16 and older to drive the cars with a moped license. The German occupation of Italy during World resulted in a severe shortage of mechanized farm equipment. The destruction of tractors was a sort of scorched-earth strategy used to reduce the independence of the conquered. The shortage of tractors in that area of Europe was the origin of Lamborghini. The war was also the inspiration for dual-purpose vehicles such as the Land Rover. Based on the Jeep, the company made a vehicle that combined PTO, tillage, 4wd, and transportation. In March 1975, a similar type of vehicle was introduced in Sweden, the A tractor [from arbetstraktor (work tractor)]; the main difference is an A tractor has a top speed of 30 km/h. This is usually done by fitting two gearboxes in a row and only using one. The Volvo Duett was, for a long time, the primary choice for conversion to an EPA or A tractor, but since supplies have dried up, other cars have been used, in most cases another Volvo. The SFRO is a Swedish organization advocating homebuilt and modified vehicles. Another type of homemade tractors are ones that are fabricated from scratch. The "from scratch" description is relative, as often individual components will be repurposed from earlier vehicles or machinery (e.g., engines, gearboxes, axle housings), but the tractor's overall chassis is essentially designed and built by the owner (e.g., a frame is welded from bar stockchannel stock, angle stock, flat stock, etc.). As with automobile conversions, the heyday of this type of tractor, at least in developed economies, lies in the past, when there were large populations of blue-collar workers for whom metalworking and farming were prevalent parts of their lives. (For example, many 19th- and 20th-century New England and Midwestern machinists and factory workers had grown up on farms.) Backyard fabrication was a natural activity to them (whereas it might seem daunting to most people today). Nomenclature The term "tractor" (US and Canada) or "tractor unit" (UK) is also applied to: Road tractors, tractor units or traction heads, familiar as the front end of an articulated lorry / semi-trailer truck. They are heavy-duty vehicles with large engines and several axles. The majority of these tractors are designed to pull long semi-trailers, most often to transport freight over a significant distance, and is connected to the trailer with a fifth wheel coupling. In England, this type of "tractor" is often called an "artic cab" (short for "articulated" cab). A minority is the ballast tractor, whose load is hauled from a drawbar. Pushback tractors are used on airports to move aircraft on the ground, most commonly pushing aircraft away from their parking stands. Locomotive tractors (engines) or rail car movers – the amalgamation of machines, electrical generators, controls and devices that comprise the traction component of railway vehicles Artillery tractors – vehicles used to tow artillery pieces of varying weights. NASA and other space agencies use very large tractors to move large launch vehicles and Space Shuttles between their hangars and launch pads. A pipe-tractor is a device used for conveying advanced instruments into pipes for measurement and data logging, and the purging of well holes, sewer pipes and other inaccessible tubes. Nebraska tests Nebraska tractor tests are tests mandated by the Nebraska Tractor Test Law and administered by the University of Nebraska, that objectively test the performance of all brands of tractors, 40 horsepower or more, sold in Nebraska. In the 1910s and 1920s, an era of snake oil sales and advertising tactics, the Nebraska tests helped farmers throughout North America to see through marketing claims and make informed buying decisions. The tests continue today, making sure tractors fulfill the manufacturer's advertised claims. Manufacturers Some of the many tractor manufacturers and brands worldwide include: Belarus Case IH Caterpillar Claas Challenger Deutz-Fahr Fendt ITMCO Iseki JCB John Deere Lamborghini Landini Kubota Mahindra Tractors Massey Ferguson McCormick Mercedes-Benz New Holland SAME Steyr TAFE Ursus Valtra Zetor In addition to commercial manufacturers, the Open Source Ecology group has developed several working prototypes of an open source hardware tractor called the LifeTrac as part of its Global Village Construction Set. Gallery	Technology	Other	null
152710	https://en.wikipedia.org/wiki/Sculptor%20%28constellation%29	Sculptor (constellation)	Sculptor is a faint constellation in the southern sky. It represents a sculptor. It was introduced by Nicolas Louis de Lacaille in the 18th century. He originally named it Apparatus Sculptoris (the sculptor's studio), but the name was later shortened. History The region to the south of Cetus and Aquarius had been named by Aratus in 270 BC as The Waters – an area of scattered faint stars with two brighter stars standing out. Professor of astronomy Bradley Schaefer has proposed that these stars were most likely Alpha and Delta Sculptoris. The French astronomer Nicolas-Louis de Lacaille first described the constellation in French as l'Atelier du Sculpteur (the sculptor's studio) in 1751–52, depicting a three-legged table with a carved head on it, and an artist's mallet and two chisels on a block of marble alongside it. Lacaille had observed and catalogued almost 10,000 southern stars during a two-year stay at the Cape of Good Hope, devising fourteen new constellations in uncharted regions of the Southern Celestial Hemisphere not visible from Europe. He named all but one in honour of instruments that symbolised the Age of Enlightenment. Characteristics Sculptor is a small constellation bordered by Aquarius and Cetus to the north, Fornax to the east, Phoenix to the south, Grus to the southwest, and Piscis Austrinus to the west. The bright star Fomalhaut is nearby. The three-letter abbreviation for the constellation, as adopted by the International Astronomical Union in 1922, is "Scl". The official constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of 6 segments. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between −24.80° and −39.37°. The whole constellation is visible to observers south of latitude 50°N. Notable features Stars No stars brighter than 3rd magnitude are located in Sculptor. This is explained by the fact that Sculptor contains the south galactic pole where stellar density is very low. Overall, there are 56 stars within the constellation's borders brighter than or equal to apparent magnitude 6.5. The brightest star is Alpha Sculptoris, an SX Arietis-type variable star with a spectral type B7IIIp and an apparent magnitude of 4.3. It is 780 ± 30 light-years distant from Earth. Eta Sculptoris is a red giant of spectral type M4III that varies between magnitudes 4.8 and 4.9, pulsating with multiple periods of 22.7, 23.5, 24.6, 47.3, 128.7 and 158.7 days. Estimated to be around 1,082 times as luminous as the Sun, it is 460 ± 20 light-years distant from Earth. R Sculptoris is a red giant that has been found to be surrounded by spirals of matter likely ejected around 1800 years ago. It is 1,440 ± 90 light-years distant from Earth. The Astronomical Society of Southern Africa in 2003 reported that observations of the Mira variable stars T, U, V and X Sculptoris were very urgently needed as data on their light curves was incomplete. Deep sky objects The constellation also contains the Sculptor Dwarf, a dwarf galaxy which is a member of the Local Group, as well as the Sculptor Group, the group of galaxies closest to the Local Group. The Sculptor Galaxy (NGC 253), a barred spiral galaxy and the largest member of the group, lies near the border between Sculptor and Cetus. Another prominent member of the group is the irregular galaxy NGC 55. One unique galaxy in Sculptor is the Cartwheel Galaxy, at a distance of 500 million light-years. The result of a merger around 300 million years ago, the Cartwheel Galaxy has a core of older, yellow stars, and an outer ring of younger, blue stars, which has a diameter of 100,000 light-years. The smaller galaxy in the collision is now incorporated into the core, after moving from a distance of 250,000 light-years. The shock waves from the collision sparked extensive star formation in the outer ring. Namesakes Sculptor (AK-103) was a United States Navy Crater class cargo ship named after the constellation.	Physical sciences	Other	Astronomy
152772	https://en.wikipedia.org/wiki/Intensive%20farming	Intensive farming	Intensive agriculture, also known as intensive farming (as opposed to extensive farming), conventional, or industrial agriculture, is a type of agriculture, both of crop plants and of animals, with higher levels of input and output per unit of agricultural land area. It is characterized by a low fallow ratio, higher use of inputs such as capital, labour, agrochemicals and water, and higher crop yields per unit land area. Most commercial agriculture is intensive in one or more ways. Forms that rely heavily on industrial methods are often called industrial agriculture, which is characterized by technologies designed to increase yield. Techniques include planting multiple crops per year, reducing the frequency of fallow years, improving cultivars, mechanised agriculture, controlled by increased and more detailed analysis of growing conditions, including weather, soil, water, weeds, and pests. Modern methods frequently involve increased use of non-biotic inputs, such as fertilizers, plant growth regulators, pesticides, and antibiotics for livestock. Intensive farms are widespread in developed nations and increasingly prevalent worldwide. Most of the meat, dairy products, eggs, fruits, and vegetables available in supermarkets are produced by such farms. Some intensive farms can use sustainable methods, although this typically necessitates higher inputs of labor or lower yields. Sustainably increasing agricultural productivity, especially on smallholdings, is an important way to decrease the amount of land needed for farming and slow and reverse environmental degradation caused by processes such as deforestation. Intensive animal farming involves large numbers of animals raised on a relatively small area of land, for example by rotational grazing, or sometimes as concentrated animal feeding operations. These methods increase the yields of food and fiber per unit land area compared to those of extensive animal husbandry; concentrated feed is brought to seldom-moved animals, or, with rotational grazing, the animals are repeatedly moved to fresh forage. History Agricultural development in Britain between the 16th century and the mid-19th century saw a massive increase in agricultural productivity and net output. This in turn contributed to unprecedented population growth, freeing up a significant percentage of the workforce, and thereby helped enable the Industrial Revolution. Historians cited enclosure, mechanization, four-field crop rotation, and selective breeding as the most important innovations. Industrial agriculture arose in the Industrial Revolution. By the early 19th century, agricultural techniques, implements, seed stocks, and cultivars had so improved that yield per land unit was many times that seen in the Middle Ages. The first phase involved a continuing process of mechanization. Horse-drawn machinery such as the McCormick reaper revolutionized harvesting, while inventions such as the cotton gin reduced the cost of processing. During this same period, farmers began to use steam-powered threshers and tractors. In 1892, the first gasoline-powered tractor was successfully developed, and in 1923, the International Harvester Farmall tractor became the first all-purpose tractor, marking an inflection point in the replacement of draft animals with machines. Mechanical harvesters (combines), planters, transplanters, and other equipment were then developed, further revolutionizing agriculture. These inventions increased yields and allowed individual farmers to manage increasingly large farms. The identification of nitrogen, phosphorus, and potassium (NPK) as critical factors in plant growth led to the manufacture of synthetic fertilizers, further increasing crop yields. In 1909, the Haber-Bosch method to synthesize ammonium nitrate was first demonstrated. NPK fertilizers stimulated the first concerns about industrial agriculture, due to concerns that they came with side effects such as soil compaction, soil erosion, and declines in overall soil fertility, along with health concerns about toxic chemicals entering the food supply. The discovery of vitamins and their role in nutrition, in the first two decades of the 20th century, led to vitamin supplements, which in the 1920s allowed some livestock to be raised indoors, reducing their exposure to adverse natural elements. Following World War II synthetic fertilizer use increased rapidly. The discovery of antibiotics and vaccines facilitated raising livestock by reducing diseases. Developments in logistics and refrigeration as well as processing technology made long-distance distribution feasible. Integrated pest management is the modern method to minimize pesticide use to more sustainable levels. There are concerns over the sustainability of industrial agriculture, and the environmental effects of fertilizers and pesticides, which has given rise to the organic movement and has built a market for sustainable intensive farming, as well as funding for the development of appropriate technology. Techniques and technologies Livestock Pasture intensification Pasture intensification is the improvement of pasture soils and grasses to increase the food production potential of livestock systems. It is commonly used to reverse pasture degradation, a process characterized by loss of forage and decreased animal carrying capacity which results from overgrazing, poor nutrient management, and lack of soil conservation. This degradation leads to poor pasture soils with decreased fertility and water availability and increased rates of erosion, compaction, and acidification. Degraded pastures have significantly lower productivity and higher carbon footprints compared to intensified pastures. Management practices which improve soil health and consequently grass productivity include irrigation, soil scarification, and the application of lime, fertilizers, and pesticides. Depending on the productivity goals of the target agricultural system, more involved restoration projects can be undertaken to replace invasive and under-productive grasses with grass species that are better suited to the soil and climate conditions of the region. These intensified grass systems allow higher stocking rates with faster animal weight gain and reduced time to slaughter, resulting in more productive, carbon-efficient livestock systems. Another technique to optimize yield while maintaining the carbon balance is the use of integrated crop-livestock (ICL) and crop-livestock-forestry (ICLF) systems, which combine several ecosystems into one optimized agricultural framework. Correctly performed, such production systems are able to create synergies potentially providing benefits to pastures through optimal plant usage, improved feed and fattening rates, increased soil fertility and quality, intensified nutrient cycling, integrated pest control, and improved biodiversity. The introduction of certain legume crops to pastures can increase carbon accumulation and nitrogen fixation in soils, while their digestibility helps animal fattening and reduces methane emissions from enteric fermentation. ICLF systems yield beef cattle productivity up to ten times that of degraded pastures; additional crop production from maize, sorghum, and soybean harvests; and greatly reduced greenhouse gas balances due to forest carbon sequestration. In the Twelve Aprils grazing program for dairy production, developed by the USDA-SARE, forage crops for dairy herds are planted into a perennial pasture. Rotational grazing Rotational grazing is a variety of foraging in which herds or flocks are regularly and systematically moved to fresh, rested grazing areas (sometimes called paddocks) to maximize the quality and quantity of forage growth. It can be used with cattle, sheep, goats, pigs, chickens, turkeys, ducks, and other animals. The herds graze one portion of pasture, or a paddock, while allowing the others to recover. Resting grazed lands allows the vegetation to renew energy reserves, rebuild shoot systems, and deepen root systems, resulting in long-term maximum biomass production. Pasture systems alone can allow grazers to meet their energy requirements, but rotational grazing is especially effective because grazers thrive on the more tender younger plant stems. Parasites are also left behind to die off, minimizing or eliminating the need for de-wormers. With the increased productivity of rotational systems, the animals may need less supplemental feed than in continuous grazing systems. Farmers can therefore increase stocking rates. Concentrated animal feeding operations Intensive livestock farming or "factory farming", is the process of raising livestock in confinement at high stocking density. "Concentrated animal feeding operations" (CAFO), or "intensive livestock operations", can hold large numbers (some up to hundreds of thousands) of cows, hogs, turkeys, or chickens, often indoors. The essence of such farms is the concentration of livestock in a given space. The aim is to provide maximum output at the lowest possible cost and with the greatest level of food safety. The term is often used pejoratively. CAFOs have dramatically increased the production of food from animal husbandry worldwide, both in terms of total food produced and efficiency. Food and water is delivered to the animals, and therapeutic use of antimicrobial agents, vitamin supplements, and growth hormones are often employed. Growth hormones are not used on chickens nor on any animal in the European Union. Undesirable behaviors often related to the stress of confinement led to a search for docile breeds (e.g., with natural dominant behaviors bred out), physical restraints to stop interaction, such as individual cages for chickens, or physical modification such as the debeaking of chickens to reduce the harm of fighting. The CAFO designation resulted from the 1972 U.S. Federal Clean Water Act, which was enacted to protect and restore lakes and rivers to a "fishable, swimmable" quality. The United States Environmental Protection Agency identified certain animal feeding operations, along with many other types of industry, as "point source" groundwater polluters. These operations were subjected to regulation. In 17 states in the U.S., isolated cases of groundwater contamination were linked to CAFOs. The U.S. federal government acknowledges the waste disposal issue and requires that animal waste be stored in lagoons. These lagoons can be as large as . Lagoons not protected with an impermeable liner can leak into groundwater under some conditions, as can runoff from manure used as fertilizer. A lagoon that burst in 1995 released 25 million gallons of nitrous sludge in North Carolina's New River. The spill allegedly killed eight to ten million fish. The large concentration of animals, animal waste, and dead animals in a small space poses ethical issues to some consumers. Animal rights and animal welfare activists have charged that intensive animal rearing is cruel to animals. Crops The Green Revolution transformed farming in many developing countries. It spread technologies that had already existed, but had not been widely used outside of industrialized nations. These technologies included "miracle seeds", pesticides, irrigation, and synthetic nitrogen fertilizer. Seeds In the 1970s, scientists created high-yielding varieties of maize, wheat, and rice. These have an increased nitrogen-absorbing potential compared to other varieties. Since cereals that absorbed extra nitrogen would typically lodge (fall over) before harvest, semi-dwarfing genes were bred into their genomes. Norin 10 wheat, a variety developed by Orville Vogel from Japanese dwarf wheat varieties, was instrumental in developing wheat cultivars. IR8, the first widely implemented high-yielding rice to be developed by the International Rice Research Institute, was created through a cross between an Indonesian variety named "Peta" and a Chinese variety named "Dee Geo Woo Gen". With the availability of molecular genetics in Arabidopsis and rice the mutant genes responsible (reduced height (rht), gibberellin insensitive (gai1) and slender rice (slr1)) have been cloned and identified as cellular signalling components of gibberellic acid, a phytohormone involved in regulating stem growth via its effect on cell division. Photosynthate investment in the stem is reduced dramatically in shorter plants and nutrients become redirected to grain production, amplifying in particular the yield effect of chemical fertilizers. High-yielding varieties outperformed traditional varieties several fold and responded better to the addition of irrigation, pesticides, and fertilizers. Hybrid vigour is utilized in many important crops to greatly increase yields for farmers. However, the advantage is lost for the progeny of the F1 hybrids, meaning seeds for annual crops need to be purchased every season, thus increasing costs and profits for farmers. Crop rotation Crop rotation or crop sequencing is the practice of growing a series of dissimilar types of crops in the same space in sequential seasons for benefits such as avoiding pathogen and pest buildup that occurs when one species is continuously cropped. Crop rotation also seeks to balance the nutrient demands of various crops to avoid soil nutrient depletion. A traditional component of crop rotation is the replenishment of nitrogen through the use of legumes and green manure in sequence with cereals and other crops. Crop rotation can also improve soil structure and fertility by alternating deep-rooted and shallow-rooted plants. A related technique is to plant multi-species cover crops between commercial crops. This combines the advantages of intensive farming with continuous cover and polyculture. Irrigation Crop irrigation accounts for 70% of the world's fresh water use. Flood irrigation, the oldest and most common type, is typically unevenly distributed, as parts of a field may receive excess water in order to deliver sufficient quantities to other parts. Overhead irrigation, using center-pivot or lateral-moving sprinklers, gives a much more equal and controlled distribution pattern. Drip irrigation is the most expensive and least-used type, but delivers water to plant roots with minimal losses. Water catchment management measures include recharge pits, which capture rainwater and runoff and use it to recharge groundwater supplies. This helps in the replenishment of groundwater wells and eventually reduces soil erosion. Dammed rivers creating reservoirs store water for irrigation and other uses over large areas. Smaller areas sometimes use irrigation ponds or groundwater. Weed control In agriculture, systematic weed management is usually required, often performed by machines such as cultivators or liquid herbicide sprayers. Herbicides kill specific targets while leaving the crop relatively unharmed. Some of these act by interfering with the growth of the weed and are often based on plant hormones. Weed control through herbicide is made more difficult when the weeds become resistant to the herbicide. Solutions include: Cover crops (especially those with allelopathic properties) that out-compete weeds or inhibit their regeneration Multiple herbicides, in combination or in rotation Strains genetically engineered for herbicide tolerance Locally adapted strains that tolerate or out-compete weeds Tilling Ground cover such as mulch or plastic Manual removal Mowing Grazing Burning Terracing In agriculture, a terrace is a leveled section of a hilly cultivated area, designed as a method of soil conservation to slow or prevent the rapid surface runoff of irrigation water. Often such land is formed into multiple terraces, giving a stepped appearance. The human landscapes of rice cultivation in terraces that follow the natural contours of the escarpments, like contour ploughing, are a classic feature of the island of Bali and the Banaue Rice Terraces in Banaue, Ifugao, Philippines. In Peru, the Inca made use of otherwise unusable slopes by building drystone walls to create terraces known as Andéns. Rice paddies A paddy field is a flooded parcel of arable land used for growing rice and other semiaquatic crops. Paddy fields are a typical feature of rice-growing countries of east and southeast Asia, including Malaysia, China, Sri Lanka, Myanmar, Thailand, Korea, Japan, Vietnam, Taiwan, Indonesia, India, and the Philippines. They are also found in other rice-growing regions such as Piedmont (Italy), the Camargue (France), and the Artibonite Valley (Haiti). They can occur naturally along rivers or marshes, or can be constructed, even on hillsides. They require large water quantities for irrigation, much of it from flooding. It gives an environment favourable to the strain of rice being grown, and is hostile to many species of weeds. As the only draft animal species which is comfortable in wetlands, the water buffalo is in widespread use in Asian rice paddies. A recent development in the intensive production of rice is the System of Rice Intensification. Developed in 1983 by the French Jesuit Father Henri de Laulanié in Madagascar, by 2013 the number of smallholder farmers using the system had grown to between 4 and 5 million. Aquaculture Aquaculture is the cultivation of the natural products of water (fish, shellfish, algae, seaweed, and other aquatic organisms). Intensive aquaculture takes place on land using tanks, ponds, or other controlled systems, or in the ocean, using cages. Sustainability Intensive farming practices which are thought to be sustainable have been developed to slow the deterioration of agricultural land and even regenerate soil health and ecosystem services. These developments may fall in the category of organic farming, or the integration of organic and conventional agriculture. Pasture cropping involves planting grain crops directly into grassland without first applying herbicides. The perennial grasses form a living mulch understory to the grain crop, eliminating the need to plant cover crops after harvest. The pasture is intensively grazed both before and after grain production. This intensive system yields equivalent farmer profits (partly from increased livestock forage) while building new topsoil and sequestering up to 33 tons of CO2/ha/year. Biointensive agriculture focuses on maximizing efficiency such as per unit area, energy input and water input. Agroforestry combines agriculture and orchard/forestry technologies to create more integrated, diverse, productive, profitable, healthy and sustainable land-use systems. Intercropping can increase yields or reduce inputs and thus represents (potentially sustainable) agricultural intensification. However, while total yield per unit land area is often increased, yields of any single crop often decrease. There are also challenges to farmers who rely on farming equipment optimized for monoculture, often resulting in increased labor inputs. Vertical farming is intensive crop production on a large scale in urban centers, in multi-story, artificially-lit structures, for the production of low-calorie foods like herbs, microgreens, and lettuce. An integrated farming system is a progressive, sustainable agriculture system such as zero waste agriculture or integrated multi-trophic aquaculture, which involves the interactions of multiple species. Elements of this integration can include: Intentionally introducing flowering plants into agricultural ecosystems to increase pollen-and nectar-resources required by natural enemies of insect pests Using crop rotation and cover crops to suppress nematodes in potatoes Integrated multi-trophic aquaculture is a practice in which the by-products (wastes) from one species are recycled to become inputs (fertilizers, food) for another. Challenges Environmental impact Industrial agriculture uses huge amounts of water, energy, and industrial chemicals, increasing pollution in the arable land, usable water, and atmosphere. Herbicides, insecticides, and fertilizers accumulate in ground and surface waters. Industrial agricultural practices are one of the main drivers of global warming, accounting for 14–28% of net greenhouse gas emissions. Many of the negative effects of industrial agriculture may emerge at some distance from fields and farms. Nitrogen compounds from the Midwest, for example, travel down the Mississippi to degrade coastal fisheries in the Gulf of Mexico, causing so-called oceanic dead zones. Many wild plant and animal species have become extinct on a regional or national scale, and the functioning of agro-ecosystems has been profoundly altered. Agricultural intensification includes a variety of factors, including the loss of landscape elements, increased farm and field sizes, and increase usage of insecticides and herbicides. The large scale of insecticides and herbicides lead to the rapid developing resistance among pests renders herbicides and insecticides increasingly ineffective. Agrochemicals have may be involved in colony collapse disorder, in which the individual members of bee colonies disappear. (Agricultural production is highly dependent on bees to pollinate many varieties of fruits and vegetables.) Intensive farming creates conditions for parasite growth and transmission that are vastly different from what parasites encounter in natural host populations, potentially altering selection on a variety of traits such as life-history traits and virulence. Some recent epidemic outbreaks have highlighted the association with intensive agricultural farming practices. For example the infectious salmon anaemia (ISA) virus is causing significant economic loss for salmon farms. The ISA virus is an orthomyxovirus with two distinct clades, one European and one North American, that diverged before 1900 (Krossøy et al. 2001). This divergence suggests that an ancestral form of the virus was present in wild salmonids prior to the introduction of cage-cultured salmonids. As the virus spread from vertical transmission (parent to offspring). Intensive monoculture increases the risk of failures due to pests, adverse weather and disease. Social impact A study for the U.S. Office of Technology Assessment concluded that regarding industrial agriculture, there is a "negative relationship between the trend toward increasing farm size and the social conditions in rural communities" on a "statistical level". Agricultural monoculture can entail social and economic risks.	Technology	Forms	null
152776	https://en.wikipedia.org/wiki/Organ%20%28biology%29	Organ (biology)	In a multicellular organism, an organ is a collection of tissues joined in a structural unit to serve a common function. In the hierarchy of life, an organ lies between tissue and an organ system. Tissues are formed from same type cells to act together in a function. Tissues of different types combine to form an organ which has a specific function. The intestinal wall for example is formed by epithelial tissue and smooth muscle tissue. Two or more organs working together in the execution of a specific body function form an organ system, also called a biological system or body system. An organ's tissues can be broadly categorized as parenchyma, the functional tissue, and stroma, the structural tissue with supportive, connective, or ancillary functions. For example, the gland's tissue that makes the hormones is the parenchyma, whereas the stroma includes the nerves that innervate the parenchyma, the blood vessels that oxygenate and nourish it and carry away its metabolic wastes, and the connective tissues that provide a suitable place for it to be situated and anchored. The main tissues that make up an organ tend to have common embryologic origins, such as arising from the same germ layer. Organs exist in most multicellular organisms. In single-celled organisms such as members of the eukaryotes, the functional analogue of an organ is known as an organelle. In plants, there are three main organs. The number of organs in any organism depends on the definition used. There are approxiamately 79 Organs in the human body,but it is something that is debated as not all scientist agree on what counts as an organ. Animals Except for placozoans, multicellular animals including humans have a variety of organ systems. These specific systems are widely studied in human anatomy. The functions of these organ systems often share significant overlap. For instance, the nervous and endocrine system both operate via a shared organ, the hypothalamus. For this reason, the two systems are combined and studied as the neuroendocrine system. The same is true for the musculoskeletal system because of the relationship between the muscular and skeletal systems. Cardiovascular system: pumping and channeling blood to and from the body and lungs with heart, blood and blood vessels. Digestive system: digestion and processing food with salivary glands, esophagus, stomach, liver, gallbladder, pancreas, intestines, colon, mesentery, rectum and anus. Endocrine system: communication within the body using hormones made by endocrine glands such as the hypothalamus, pituitary gland, pineal body or pineal gland, thyroid, parathyroids and adrenals, i.e., adrenal glands. Excretory system: kidneys, ureters, bladder and urethra involved in fluid balance, electrolyte balance and excretion of urine. Lymphatic system: structures involved in the transfer of lymph between tissues and the blood stream, the lymph and the nodes and vessels that transport it including the immune system: defending against disease-causing agents with leukocytes, tonsils, adenoids, thymus and spleen. Integumentary system: skin, hair and nails of mammals. Also scales of fish, reptiles, and birds, and feathers of birds. Muscular system: movement with muscles. Nervous system: collecting, transferring and processing information with brain, spinal cord and nerves. Reproductive system: the sex organs, such as ovaries, oviducts, uterus, vulva, vagina, testicles, vasa deferentia, seminal vesicles, prostate and penis. Respiratory system: the organs used for breathing, the pharynx, larynx, trachea, bronchi, lungs and diaphragm. Skeletal system: structural support and protection with bones, cartilage, ligaments and tendons. Viscera In the study of anatomy, viscera (: viscus) refers to the internal organs of the abdominal, thoracic, and pelvic cavities. The abdominal organs may be classified as solid organs or hollow organs. The solid organs are the liver, pancreas, spleen, kidneys, and adrenal glands. The hollow organs of the abdomen are the stomach, intestines, gallbladder, bladder, and rectum. In the thoracic cavity, the heart is a hollow, muscular organ. Splanchnology is the study of the viscera. The term "visceral" is contrasted with the term "", meaning "of or relating to the wall of a body part, organ or cavity". The two terms are often used in describing a membrane or piece of connective tissue, referring to the opposing sides. Origin and evolution The organ level of organisation in animals can be first detected in flatworms and the more derived phyla, i.e. the bilaterians. The less-advanced taxa (i.e. Placozoa, Porifera, Ctenophora and Cnidaria) do not show unification of their tissues into organs. More complex animals are composed of different organs, which have evolved over time. For example, the liver and heart evolved in the chordates about 550-500 million years ago, while the gut and brain are even more ancient, arising in the ancestor of vertebrates, insects, molluscs, and worms about 700–650 million years ago. Given the ancient origin of most vertebrate organs, researchers have looked for model systems, where organs have evolved more recently, and ideally have evolved multiple times independently. An outstanding model for this kind of research is the placenta, which has evolved more than 100 times independently in vertebrates, has evolved relatively recently in some lineages, and exists in intermediate forms in extant taxa. Studies on the evolution of the placenta have identified a variety of genetic and physiological processes that contribute to the origin and evolution of organs, these include the re-purposing of existing animal tissues, the acquisition of new functional properties by these tissues, and novel interactions of distinct tissue types. Plants The study of plant organs is covered in plant morphology. Organs of plants can be divided into vegetative and reproductive. Vegetative plant organs include roots, stems, and leaves. The reproductive organs are variable. In flowering plants, they are represented by the flower, seed and fruit. In conifers, the organ that bears the reproductive structures is called a cone. In other divisions (phyla) of plants, the reproductive organs are called strobili, in Lycopodiophyta, or simply gametophores in mosses. Common organ system designations in plants include the differentiation of shoot and root. All parts of the plant above ground (in non-epiphytes), including the functionally distinct leaf and flower organs, may be classified together as the shoot organ system. The vegetative organs are essential for maintaining the life of a plant. While there can be 11 organ systems in animals, there are far fewer in plants, where some perform the vital functions, such as photosynthesis, while the reproductive organs are essential in reproduction. However, if there is asexual vegetative reproduction, the vegetative organs are those that create the new generation of plants (see clonal colony). Society and culture Many societies have a system for organ donation, in which a living or deceased donor's organ are transplanted into a person with a failing organ. The transplantation of larger solid organs often requires immunosuppression to prevent organ rejection or graft-versus-host disease. There is considerable interest throughout the world in creating laboratory-grown or artificial organs. Organ transplants Beginning in the 20th century, organ transplants began to take place as scientists knew more about the anatomy of organs. These came later in time as procedures were often dangerous and difficult. Both the source and method of obtaining the organ to transplant are major ethical issues to consider, and because organs as resources for transplant are always more limited than demand for them, various notions of justice, including distributive justice, are developed in the ethical analysis. This situation continues as long as transplantation relies upon organ donors rather than technological innovation, testing, and industrial manufacturing. History The English word "organ" dates back to the twelfth century and refers to any musical instrument. By the late 14th century, the musical term's meaning had narrowed to refer specifically to the keyboard-based instrument. At the same time, a second meaning arose, in reference to a "body part adapted to a certain function". Plant organs are made from tissue composed of different types of tissue. The three tissue types are ground, vascular, and dermal. When three or more organs are present, it is called an organ system. The adjective visceral, also splanchnic, is used for anything pertaining to the internal organs. Historically, viscera of animals were examined by Roman pagan priests like the haruspices or the augurs in order to divine the future by their shape, dimensions or other factors. This practice remains an important ritual in some remote, tribal societies. The term "visceral" is contrasted with the term "", meaning "of or relating to the wall of a body part, organ or cavity" The two terms are often used in describing a membrane or piece of connective tissue, referring to the opposing sides. Antiquity Aristotle used the word frequently in his philosophy, both to describe the organs of plants or animals (e.g. the roots of a tree, the heart or liver of an animal) because, in ancient Greek, the word 'organon' means 'tool', and Aristotle believed that the organs of the body were tools for us by means of which we can do things. For similar reasons, his logical works, taken as a whole, are referred to as the Organon because logic is a tool for philosophical thinking. Earlier thinkers, such as those who wrote texts in the Hippocratic corpus, generally did not believe that there were organs of the body but only different parts of the body. Some alchemists (e.g. Paracelsus) adopted the Hermetic Qabalah assignment between the seven vital organs and the seven classical planets as follows: Chinese traditional medicine recognizes eleven organs, associated with the five Chinese traditional elements and with yin and yang, as follows: The Chinese associated the five elements with the five planets (Jupiter, Mars, Venus, Saturn, and Mercury) similar to the way the classical planets were associated with different metals. The yin and yang distinction approximates the modern notion of solid and hollow organs.	Biology and health sciences	Basics_2	null
152952	https://en.wikipedia.org/wiki/Portuguese%20man%20o%27%20war	Portuguese man o' war	The Portuguese war (Physalia physalis), also known as the man-of-war or bluebottle, is a marine hydrozoan found in the Atlantic Ocean and the Indian Ocean. It is considered to be the same species as the Pacific man o' war or bluebottle, which is found mainly in the Pacific Ocean. The Portuguese man o' war is the only species in the genus Physalia, which in turn is the only genus in the family Physaliidae. The Portuguese man o' war is a conspicuous member of the neuston, the community of organisms that live at the surface of the ocean. It has numerous microscopic venomous cnidocytes which deliver a painful sting powerful enough to kill fish, and even, in some cases, humans. Although it superficially resembles a jellyfish, the Portuguese man o' war is in fact a siphonophore. Like all siphonophores, it is a colonial organism, made up of many smaller units called zooids. Although they are morphologically quite different, all of the zooids in a single specimen are genetically identical. These different types of zooids fulfill specialized functions, such as hunting, digestion and reproduction, and together they allow the colony to operate as a single individual. Etymology The name man o’ war comes from the man-of-war, a sailing warship, and the animal's resemblance to the Portuguese version (the caravel) at full sail. Taxonomy The bluebottle, Pacific man o' war or Indo-Pacific Portuguese man o' war, distinguished by a smaller float and a single long fishing tentacle, was originally considered a separate species in the same genus (P. utriculus). The name was synonymized with P. physalis in 2007, and it is now considered a regional form of the same species. Coloniality The man o' war is described as a colonial organism because the individual zooids in a colony are evolutionarily derived from either polyps or medusae, i.e. the two basic body plans of cnidarians. Both of these body plans comprise entire individuals in non-colonial cnidarians (for example, a jellyfish is a medusa, while a sea anemone is a polyp). All zooids in a man o' war develop from the same single fertilized egg and are therefore genetically identical. They remain physiologically connected throughout life, and essentially function as organs in a shared body. Hence, a Portuguese man o' war constitutes a single organism from an ecological perspective, but is made up of many individuals from an embryological perspective. Most species of siphonophores are fragile and difficult to collect intact. However, P. physalis is the most accessible, conspicuous, and robust of the siphonophores, and much has been written about this species. The development, morphology, and colony organization of P. physalis is very different from that of other siphonophores. Its structure, embryological development, and histology have been examined by several authors. These studies provide an important foundation for understanding the morphology, cellular anatomy, and development of this species. Description Like all siphonophores, P. physalis is a colonial organism: each animal is composed of many smaller units (zooids) that hang in clusters from under a large, gas-filled structure called the pneumatophore. Seven different types of zooids have been described in the man o' war, and all of these are interdependent on each other for survival and performing different functions, such as digestion (gastrozooids), reproduction (gonozooids) and hunting (dactylozooids). A fourth type of zooid is the pneumatophore. Three of these types of zooids are of the medusoid type (gonophores, nectophores, and vestigial nectophores), while the remaining four are of the polypoid type (free gastrozooids, tentacle-bearing zooids, gonozooids and gonopalpons). However, naming and categorization of zooids varies between authors, and much of the embryonic and evolutionary relationships of zooids remains unclear. The pneumatophore or bladder is the most conspicuous part of the man o' war. This large, gas-filled, translucent structure is pink, purple or blue in color; it is long and rises as much as above the water. The pneumatophore functions as both a flotation device and a sail, allowing the animal to move with the prevailing wind. The gas in the pneumatophore is mostly air which diffuses in from the surrounding atmosphere, but it also contains as much as 13% carbon monoxide, which is actively produced by the animal. In the event of a surface attack, the pneumatophore can be deflated, allowing the animal to temporarily submerge. New zooids are added by budding as the colony grows. Long tentacles hang below the float as the animal drifts, fishing for prey to sting and drag up to its digestive zooids. The colony hunts and feeds through the cooperation of two types of zooids: tentacle-bearing zooids known as dactylozooids (or palpons), and gastrozooids. The palpons are equipped with tentacles, which are typically about in length but can reach over . Each tentacle bears tiny, coiled, thread-like structures called nematocysts. Nematocysts trigger and inject venom on contact, stinging, paralyzing, and killing molluscs and fishes. Large groups of Portuguese man o' war, sometimes over 1,000 individuals, may deplete fisheries. Contraction of tentacles drags the prey upward and into range of the gastrozooids. The gastrozooids surround and digest the food by secreting digestive enzymes. P. physalis typically has multiple stinging tentacles, but a regional form (previously known as a separate species, P. utriculus) has only a single stinging tentacle. The main reproductive zooids, the gonophores, are situated on branching structures called gonodendra. Gonophores produce sperm or eggs. Besides gonophores, each gonodendron also contains several other types of specialized zooids: gonozooids (which are accessory gastrozooids), nectophores (which have been speculated to allow detached gonodendra to swim), and vestigial nectophores (also called jelly polyps; the function of these is unclear). Life cycle Man o' war individuals are dioecious, meaning each colony is either male or female. Gonophores producing either sperm or eggs (depending on the sex of the colony) sit on a tree-like structure called a gonodendron, which is believed to drop off from the colony during reproduction. Mating takes place primarily in the autumn, when eggs and sperm are shed from gonophores into the water. As neither fertilization nor early development has been directly observed in the wild, it is not yet known at what depth these occur. A fertilized man o' war egg develops into a planula that buds off new zooids as it grows, gradually forming a new colony. This development initially occurs under the water, and has been reconstructed by comparing different stages of planulae collected at sea. The first two structures to emerge are the pneumatophore (sail) and a single, early feeding zooid called a protozooid. Later, gastrozooids and tentacle-bearing zooids are added. Eventually, the growing pneumatophore becomes buoyant enough to carry the immature colony on the surface of the water. Ecology Predators and prey The Portuguese man o' war is a carnivore. Using its venomous tentacles, it traps and paralyzes its prey while reeling it inwards to its digestive polyps. It typically feeds on small fish, molluscs, shrimp and other small crustaceans, and zooplankton. The organism has few predators; one example is the loggerhead sea turtle, which feeds on the Portuguese man o' war as a common part of its diet. The turtle's skin, including that of its tongue and throat, is too thick for the stings to penetrate. Also, the blue sea slug specializes in feeding on the Portuguese man o' war, as does the violet sea snail. The ocean sunfish's diet, once thought to consist mainly of jellyfish, has been found to include many species, including the Portuguese man o' war. The man-of-war fish, Nomeus gronovii, is a driftfish native to the Atlantic, Pacific and Indian Oceans. It is notable for its ability to live within the deadly tentacles of the Portuguese man o' war, upon whose tentacles and gonads it feeds. Rather than using mucus to prevent nematocysts from firing, as is seen in some of the clownfish sheltering among sea anemones, the man-of-war fish appears to use highly agile swimming to physically avoid tentacles. The fish has a very high number of vertebrae (41), which may add to its agility and primarily uses its pectoral fins for swimming—a feature of fish that specialize in maneuvering tight spaces. It also has a complex skin design and at least one antibody to the man o' war's toxins. Although the fish seems to be 10 times more resistant to the toxin than other fish, it can be stung by the dactylozooides (large tentacles), which it actively avoids. The smaller gonozooids do not seem to sting the fish and the fish is reported to frequently nibble on these tentacles. Commensalism and symbiosis The Portuguese man o' war is often found with a variety of other marine fish, including yellow jack. These fish benefit from the shelter from predators provided by the stinging tentacles, and for the Portuguese , the presence of these species may attract other fish to eat. The blanket octopus is immune to the venom of the Portuguese man o' war. Individuals have been observed to carry broken man o' war tentacles, which males and immature females rip off and use for offensive and defensive purposes. Venom The stinging, venom-filled nematocysts in the tentacles of the Portuguese man o' war can paralyze small fish and other prey. Detached tentacles and dead specimens (including those that wash up on shore) can sting just as painfully as those of the live organism in the water and may remain potent for hours or even days after the death of the organism or the detachment of the tentacle. Stings usually cause severe pain to humans, lasting one to three hours. Red, whip-like welts appear on the skin that last two or three days after the sting. In some cases, the venom may travel to the lymph nodes and may cause symptoms that mimic an allergic reaction, including swelling of the larynx, airway blockage, cardiac distress and shortness of breath. Other symptoms may include fever, circulatory shock and in extreme cases, even death, although this is extremely rare. Medical attention for those exposed to large numbers of tentacles may become necessary to relieve pain or open airways if the pain becomes excruciating or lasts for more than three hours, or if breathing becomes difficult. Instances in which the stings completely surround the trunk of a young child are among those that may be fatal. The species is responsible for up to 10,000 human stings in Australia each summer, particularly on the east coast, with some others occurring off the coast of South Australia and Western Australia. Treatment of stings Stings from a Portuguese man o' war can result in severe dermatitis characterized by long, thin, open wounds that resemble those caused by a whip. These are not caused by any impact or cutting action, but by irritating urticariogenic substances in the tentacles. Treatment for sting pain is immersion in water for 20 minutes. The cnidocyte found in the box jellyfish react differently than the nematocyst in the Portuguese man o' war; cnidocytes are inhibited by application of vinegar, but nematocysts can discharge more venom if vinegar is applied. Distribution The species is found throughout the world's oceans, mainly in tropical and subtropical regions, but occasionally also in temperate regions. Habitat P. physalis is a member of the neuston (the floating community of organisms that live at the interface between water and air). This community is exposed to a unique set of environmental conditions including prolonged exposure to intense ultraviolet light, risk of desiccation, and rough sea conditions. The gas-filled bladder, or pneumatophore, remains at the surface, while the remainder is submerged. The animal has no means of propulsion; it moves passively, driven by the winds, currents, and tides. Winds can drive them into bays or onto beaches. Often, finding a single Portuguese man o' war is followed by finding many others in the vicinity. The Portuguese man o' war is well known to beachgoers for the painful stings delivered by its tentacles. Because they can sting while beached, the discovery of a man o' war washed up on a beach may lead to the closure of the beach. Drifting dynamics P. physalis uses a float filled with carbon monoxide and air as a sail to travel by wind for thousands of miles, dragging behind long tentacles that deliver a deadly venomous sting to fish. This sailing ability, combined with a painful sting and a life cycle with seasonal blooms, results in periodic mass beach strandings and occasional human envenomations, making P. physalis the most infamous of the siphonophores. Despite being a common occurrence, the origin of the man o' war or bluebottle before reaching the coastline is not well understood, and neither is the way it drifts at the surface of the ocean. Left- and right-handedness The Portuguese man o' war is asymmetrically shaped: the zooids hang down from either the right or left side of the midline of the pneumatophore or bladder. The pneumatophore can be oriented towards the left or the right. This phenomenon may be an adaptation that prevents an entire population from being washed on shore to die. The "left-handed" animals sail to the right of the wind, while the "right-handed" animals sail to the left. The wind will always push the two types in opposite directions, so at most half the population will be pushed towards the coast. Regional populations can have substantial differences in float size and the number of tentacles used for hunting. The regional form previously known as P. utriculus has a bladder rarely exceeding in length and has one long hunting tentacle that is less than long. In comparison, the typical man o' war has a float of around , and several hunting tentacles that can reach in mature colonies when fully extended. When combined with the trailing action of the tentacles, this left- or right-handedness makes the colony sail sideways relative to the wind, by about 45° in either direction. Colony handedness has therefore been theorized to influence man o' war migration, with left-handed or right-handed colonies potentially being more likely to drift down particular respective sea routes. Handedness develops early in the colony's life, while it is still living below the surface of the sea. Mathematical modelling Since they have no propulsion system, the movement of the man o' war can be modelled mathematically by calculating the forces acting on it, or by advecting virtual particles in ocean and atmospheric circulation models. Earlier studies modelled the movement of the man o' war with Lagrangian particle tracking to explain major beaching events. In 2017, Ferrer and Pastor were able to estimate the region of origin of a significant beaching event on the southeastern Bay of Biscay. They ran a Lagrangian model backwards in time, using wind velocity and a wind drag coefficient as drivers of the man o' war motion. They found that the region of origin was the North Atlantic subtropical gyre. In 2015 Prieto et al. included both the effect of the surface currents and wind to predict the initial colony position prior to major beaching events in the Mediterranean. This model assumed the man o' war was advected by the surface currents, with the effect of the wind being added with a much higher wind drag coefficient of 10%. Similarly, in 2020 Headlam et al. used beaching and offshore observations to identify a region of origin, using the joint effects of surface currents and wind drag, for the largest mass man o' war beaching on the Irish coastline in over 150 years. These earlier studies used numerical models in combination with simple assumptions to calculate the drift of this species, excluding complex drifting dynamics. In 2021, Lee et al. provide a parameterisation for Lagrangian modelling of the bluebottle by considering the similarities between the bluebottle and a sailboat. This allowed them to compute the hydrodynamic and aerodynamic forces acting on the bluebottle and use an equilibrium condition to create a generalised model for calculating the drifting speed and course of the bluebottle under any wind and ocean current conditions. Gallery	Biology and health sciences	Cnidarians	Animals
152969	https://en.wikipedia.org/wiki/Eutectic%20system	Eutectic system	A eutectic system or eutectic mixture ( ) is a type of a homogeneous mixture that has a melting point lower than those of the constituents. The lowest possible melting point over all of the mixing ratios of the constituents is called the eutectic temperature. On a phase diagram, the eutectic temperature is seen as the eutectic point (see plot on the right). Non-eutectic mixture ratios have different melting temperatures for their different constituents, since one component's lattice will melt at a lower temperature than the other's. Conversely, as a non-eutectic mixture cools down, each of its components solidifies into a lattice at a different temperature, until the entire mass is solid. A non-eutectic mixture thus does not have a single melting/freezing point temperature at which it changes phase, but rather a temperature at which it changes between liquid and slush (known as the liquidus) and a lower temperature at which it changes between slush and solid (the solidus). In the real world, eutectic properties can be used to advantage in such processes as eutectic bonding, where silicon chips are bonded to gold-plated substrates with ultrasound, and eutectic alloys prove valuable in such diverse applications as soldering, brazing, metal casting, electrical protection, fire sprinkler systems, and nontoxic mercury substitutes. The term was coined in 1884 by British physicist and chemist Frederick Guthrie (1833–1886). The word originates . Before his studies, chemists assumed "that the alloy of minimum fusing point must have its constituents in some simple atomic proportions", which was indeed proven to be not the case. Eutectic phase transition The eutectic solidification is defined as follows: This type of reaction is an invariant reaction, because it is in thermal equilibrium; another way to define this is the change in Gibbs free energy equals zero. Tangibly, this means the liquid and two solid solutions all coexist at the same time and are in chemical equilibrium. There is also a thermal arrest for the duration of the phase change during which the temperature of the system does not change. The resulting solid macrostructure from a eutectic reaction depends on a few factors, with the most important factor being how the two solid solutions nucleate and grow. The most common structure is a lamellar structure, but other possible structures include rodlike, globular, and acicular. Non-eutectic compositions Compositions of eutectic systems that are not at the eutectic point can be classified as hypoeutectic or hypereutectic: Hypoeutectic compositions are those with a greater composition of species α and a smaller percent composition of species β than the eutectic composition (E) Hypereutectic compositions are characterized as those with a higher composition of species β and a lower composition of species α than the eutectic composition. As the temperature of a non-eutectic composition is lowered the liquid mixture will precipitate one component of the mixture before the other. In a hypereutectic solution, there will be a proeutectoid phase of species β whereas a hypoeutectic solution will have a proeutectic α phase. Types Alloys Eutectic alloys have two or more materials and have a eutectic composition. When a non-eutectic alloy solidifies, its components solidify at different temperatures, exhibiting a plastic melting range. Conversely, when a well-mixed, eutectic alloy melts, it does so at a single, sharp temperature. The various phase transformations that occur during the solidification of a particular alloy composition can be understood by drawing a vertical line from the liquid phase to the solid phase on the phase diagram for that alloy. Some uses for eutectic alloys include: NEMA eutectic alloy overload relays for electrical protection of three-phase motors for pumps, fans, conveyors, and other factory process equipment. Eutectic alloys for soldering, both traditional alloys composed of lead (Pb) and tin (Sn), sometimes with additional silver (Ag) or gold (Au) — especially SnPb and SnPbAg alloy formula for electronics - and newer lead-free soldering alloys, in particular ones composed of tin, silver, and copper (Cu) such as SnAg. Casting alloys, such as aluminium-silicon and cast iron (at the composition of 4.3% carbon in iron producing an austenite-cementite eutectic) Silicon chips are eutectic bonded to gold-plated substrates through a silicon-gold eutectic by the application of ultrasonic energy to the chip. Brazing, where diffusion can remove alloying elements from the joint, so that eutectic melting is only possible early in the brazing process Temperature response, e.g., Wood's metal and Field's metal for fire sprinklers Non-toxic mercury replacements, such as galinstan Experimental glassy metals, with extremely high strength and corrosion resistance Eutectic alloys of sodium and potassium (NaK) that are liquid at room temperature and used as coolant in experimental fast neutron nuclear reactors. Others Sodium chloride and water form a eutectic mixture whose eutectic point is −21.2 °C and 23.3% salt by mass. The eutectic nature of salt and water is exploited when salt is spread on roads to aid snow removal, or mixed with ice to produce low temperatures (for example, in traditional ice cream making). Ethanol–water has an unusually biased eutectic point, i.e. it is close to pure ethanol, which sets the maximum proof obtainable by fractional freezing. "Solar salt", 60% NaNO3 and 40% KNO3, forms a eutectic molten salt mixture which is used for thermal energy storage in concentrated solar power plants. To reduce the eutectic melting point in the solar molten salts, calcium nitrate is used in the following proportion: 42% Ca(NO3)2, 43% KNO3, and 15% NaNO3. Lidocaine and prilocaine—both are solids at room temperature—form a eutectic that is an oil with a melting point that is used in eutectic mixture of local anesthetic (EMLA) preparations. Menthol and camphor, both solids at room temperature, form a eutectic that is a liquid at room temperature in the following proportions: 8:2, 7:3, 6:4, and 5:5. Both substances are common ingredients in pharmacy extemporaneous preparations. Minerals may form eutectic mixtures in igneous rocks, giving rise to characteristic intergrowth textures exhibited, for example, by granophyre. Some inks are eutectic mixtures, allowing inkjet printers to operate at lower temperatures. Choline chloride produces eutectic mixtures with many natural products such as citric acid, malic acid and sugars. These liquid mixtures can be used, for example, to obtain antioxidant and antidiabetic extracts from natural products. Strengthening mechanisms Alloys The primary strengthening mechanism of the eutectic structure in metals is composite strengthening (See strengthening mechanisms of materials). This deformation mechanism works through load transfer between the two constituent phases where the more compliant phase transfers stress to the stiffer phase. By taking advantage of the strength of the stiff phase and the ductility of the compliant phase, the overall toughness of the material increases. As the composition is varied to either hypoeutectic or hypereutectic formations, the load transfer mechanism becomes more complex as there is a load transfer between the eutectic phase and the secondary phase as well as the load transfer within the eutectic phase itself. A second tunable strengthening mechanism of eutectic structures is the spacing of the secondary phase. By changing the spacing of the secondary phase, the fraction of contact between the two phases through shared phase boundaries is also changed. By decreasing the spacing of the eutectic phase, creating a fine eutectic structure, more surface area is shared between the two constituent phases resulting in more effective load transfer. On the micro-scale, the additional boundary area acts as a barrier to dislocations further strengthening the material. As a result of this strengthening mechanism, coarse eutectic structures tend to be less stiff but more ductile while fine eutectic structures are stiffer but more brittle. The spacing of the eutectic phase can be controlled during processing as it is directly related to the cooling rate during solidification of the eutectic structure. For example, for a simple lamellar eutectic structure, the minimal lamellae spacing is: Where is is the surface energy of the two-phase boundary, is the molar volume of the eutectic phase, is the solidification temperature of the eutectic phase, is the enthalpy of formation of the eutectic phase, and is the undercooling of the material. So, by altering the undercooling, and by extension the cooling rate, the minimal achievable spacing of the secondary phase is controlled. Strengthening metallic eutectic phases to resist deformation at high temperatures (see creep deformation) is more convoluted as the primary deformation mechanism changes depending on the level of stress applied. At high temperatures where deformation is dominated by dislocation movement, the strengthening from load transfer and secondary phase spacing remain as they continue to resist dislocation motion. At lower strains where Nabarro-Herring creep is dominant, the shape and size of the eutectic phase structure plays a significant role in material deformation as it affects the available boundary area for vacancy diffusion to occur. Other critical points Eutectoid When the solution above the transformation point is solid, rather than liquid, an analogous eutectoid transformation can occur. For instance, in the iron-carbon system, the austenite phase can undergo a eutectoid transformation to produce ferrite and cementite, often in lamellar structures such as pearlite and bainite. This eutectoid point occurs at and 0.76 wt% carbon. Peritectoid A peritectoid transformation is a type of isothermal reversible reaction that has two solid phases reacting with each other upon cooling of a binary, ternary, ..., n-ary alloy to create a completely different and single solid phase. The reaction plays a key role in the order and decomposition of quasicrystalline phases in several alloy types. A similar structural transition is also predicted for rotating columnar crystals. Peritectic Peritectic transformations are also similar to eutectic reactions. Here, a liquid and solid phase of fixed proportions react at a fixed temperature to yield a single solid phase. Since the solid product forms at the interface between the two reactants, it can form a diffusion barrier and generally causes such reactions to proceed much more slowly than eutectic or eutectoid transformations. Because of this, when a peritectic composition solidifies it does not show the lamellar structure that is found with eutectic solidification. Such a transformation exists in the iron-carbon system, as seen near the upper-left corner of the figure. It resembles an inverted eutectic, with the δ phase combining with the liquid to produce pure austenite at and 0.17% carbon. At the peritectic decomposition temperature the compound, rather than melting, decomposes into another solid compound and a liquid. The proportion of each is determined by the lever rule. In the Al-Au phase diagram, for example, it can be seen that only two of the phases melt congruently, AuAl2 and Au2Al, while the rest peritectically decompose. "Bad solid solution" Not all minimum melting point systems are "eutectic". The alternative of "poor solid solution" can be illustrated by comparing the common precious metal systems Cu-Ag and Cu-Au. Cu-Ag, source for example https://himikatus.ru/art/phase-diagr1/Ag-Cu.php, is a true eutectic system. The eutectic melting point is at 780 °C, with solid solubility limits at fineness 80 and 912 by weight, and eutectic at 719. Since Cu-Ag is a true eutectic, any silver with fineness anywhere between 80 and 912 will reach solidus line, and therefore melt at least partly, at exactly 780 °C. The eutectic alloy with fineness exactly 719 will reach liquidus line, and therefore melt entirely, at that exact temperature without any further rise of temperature till all of the alloy has melted. Any silver with fineness between 80 and 912 but not exactly 719 will also reach the solidus line at exactly 780 °C, but will melt partly. It will leave a solid residue with fineness of either exactly 912 or exactly 80, but never some of both. It will melt at constant temperature without further rise of temperature until the exact amount of eutectic (fineness 719) alloy has melted off to divide the alloy into eutectic melt and solid solution residue. On further heating, the solid solution residue dissolves in the melt and changes its composition until the liquidus line is reached and the whole residue has dissolved away. Cu-Au source for example https://himikatus.ru/art/phase-diagr1/Au-Cu.php does display a melting point minimum at 910 °C and given as 44 atom % Cu, which converts to about 20 weight percent Cu - about 800 fineness of gold. But this is not a true eutectic. 800 fine gold melts at 910 °C, to a melt of exact same composition, and the whole alloy will melt at exact same temperature. But the differences happen away from the minimum composition. Unlike silver with fineness other than 719 (which melts partly at exactly 780 °C through a wide fineness range), gold with fineness other than 800 will reach solidus and start partial melting at a temperature different from and higher than 910 °C, depending on the alloy fineness. The partial melting does cause some composition changes - the liquid will be closer in fineness towards 800 than the remaining solid, but the liquid will not have fineness of exactly 800 and the fineness of the remaining solid will depend on the fineness of the liquid. The underlying reason is that for an eutectic system like Cu-Ag, the solubility in liquid phase is good but solubility in solid phase is limited. Therefore when a silver-copper alloy is frozen, it actually separates into crystals of 912 fineness silver and 80 fineness silver - both are saturated and always have the same composition at the freezing point of 780 °C. Thus the alloy just below 780 °C consists of two types of crystals of exactly the same composition regardless of the total alloy composition, only the relative amount of each type of crystals differs. Therefore they always melt at 780 °C until one or other type of crystals, or both, will be exhausted. In contrast, in Cu-Au system the components are miscible at the melting point in all compositions even in solid. There can be crystals of any composition, which will melt at different temperatures depending on composition. However, Cu-Au system is a "poor" solid solution. There is a substantial misfit between the atoms in solid which, however, near the melting point is overcome by entropy of thermal motion mixing the atoms. That misfit, however, disfavours the Cu-Au solution relative to phases in which the atoms are better fitted, such as the melt, and causes the melting point to fall below the melting point of components. Eutectic calculation The composition and temperature of a eutectic can be calculated from enthalpy and entropy of fusion of each components. The Gibbs free energy G depends on its own differential: Thus, the G/T derivative at constant pressure is calculated by the following equation: The chemical potential is calculated if we assume that the activity is equal to the concentration: At the equilibrium, , thus is obtained as Using and integrating gives The integration constant K may be determined for a pure component with a melting temperature and an enthalpy of fusion : We obtain a relation that determines the molar fraction as a function of the temperature for each component: The mixture of n components is described by the system which can be solved by	Physical sciences	Phase separations	Chemistry
152970	https://en.wikipedia.org/wiki/Euglenid	Euglenid	Euglenids or euglenoids are one of the best-known groups of eukaryotic flagellates: single-celled organisms with flagella, or whip-like tails. They are classified in the phylum Euglenophyta, class Euglenida or Euglenoidea. Euglenids are commonly found in fresh water, especially when it is rich in organic materials, but they have a few marine and endosymbiotic members. Many euglenids feed by phagocytosis, or strictly by diffusion. A monophyletic subgroup known as Euglenophyceae have chloroplasts and produce their own food through photosynthesis. This group contains the carbohydrate paramylon. Euglenids split from other Euglenozoa (a larger group of flagellates) more than a billion years ago. The plastids (membranous organelles) in all extant photosynthetic species result from secondary endosymbiosis between a euglenid and a green alga. Structure Euglenoids are distinguished mainly by the presence of a type of cell covering called a pellicle. Within its taxon, the pellicle is one of the euglenoids' most diverse morphological features. The pellicle is composed of proteinaceous strips underneath the cell membrane, supported by dorsal and ventral microtubules. This varies from rigid to flexible, and gives the cell its shape, often giving it distinctive striations. In many euglenids, the strips can slide past one another, causing an inching motion called metaboly. Otherwise, they move using their flagella. Classification The first attempt at classifying euglenids was done by Ehrenberg in 1830, when he described the genus Euglena and placed it in the Polygastrica of family Astasiae, containing other creatures of variable body shape and lacking pseudopods or lorica. Later, various biologists described additional characteristics for Euglena and established different classification systems for euglenids based on nutrition modes, the presence and number of flagella, and the degree of metaboly. The 1942 revision by A. Hollande distinguished three groups, Peranemoidées (flexible phagotrophs), Petalomonadinées (rigid phagotrophs) and Euglenidinées (phototrophs), and was widely accepted as the best reflection of the natural relationships between euglenids, adopted by many other authors. Gordon F. Leedale expanded on Hollande's system, establishing six orders (Eutreptiales, Euglenales, Rhabdomonadales, Sphenomonadales, Heteronematales and Euglenamorphales) and taking into account new data on their physiology and ultrastructure. This scheme endured until 1986, with the sequencing of the SSU rRNA gene from Euglena gracilis. Euglenids are currently regarded as a highly diverse clade within Euglenozoa, in the eukaryotic supergroup Discoba. They are traditionally organized into three categories based on modes of nutrition: the phototrophs (Euglenophyceae), the osmotrophs (mainly the 'primary osmotrophs' known as Aphagea), and the phagotrophs, from which the first two groups have evolved. The phagotrophs, although paraphyletic, have historically been classified under the name of Heteronematina. In addition, euglenids can be divided into inflexible or rigid euglenids, and flexible or metabolic euglenids which are capable of 'metaboly' or 'euglenid motion'. Only those with more than 18 protein strips in their pellicle gain this flexibility. Phylogenetic studies show that various clades of rigid phagotrophic euglenids compose the base of the euglenid tree, namely Petalomonadida and the paraphyletic 'Ploeotiida'. In contrast, all flexible euglenids belong to a monophyletic group known as Spirocuta, which includes Euglenophyceae, Aphagea and various phagotrophs (Peranemidae, Anisonemidae and Neometanemidae). The current classification of class Euglenida, as a result of these studies, is as follows: Euglenida incertae sedis: Atraktomonas, Calycimonas, Dolium, Dylakosoma, Tropidoscyphus, Michajlowastasia, Parastasiella, Dinemula, Paradinemula, Mononema, Ovicola, Naupliicola, Embryocola, Copromonas. Order Petalomonadida Order "Ploeotiida" (paraphyletic) Clade Alistosa Entosiphon Gaulosia Clade Karavia Chelandium Olkasia Clade Spirocuta [Helicales ] Clade Anisonemia Order Anisonemida Family Anisonemidae Order Natomonadida Suborder Metanemina Family Neometanemidae Suborder Aphagea [Rhabdomonadina ] Family Astasiidae Family Distigmidae Order Peranemida Family Peranemidae Clade Euglenophyceae [Euglenea ] Euglenophyceae incertae sedis: Ascoglena, Euglenamorpha, Euglenopsis, Glenoclosteroium, Hegneria, Klebsina, Euglenocapsa. Order Rapazida Family Rapazidae Order Eutreptiales Family Eutreptiaceae Order Euglenales Family Phacaceae Family Euglenaceae Nutrition The classification of euglenids is still variable, as groups are being revised to conform with their molecular phylogeny. Classifications have fallen in line with the traditional groups based on differences in nutrition and number of flagella; these provide a starting point for considering euglenid diversity. Different characteristics of the euglenids' pellicles can provide insight into their modes of movement and nutrition. As with other Euglenozoa, the primitive mode of nutrition is phagocytosis. Prey such as bacteria and smaller flagellates is ingested through a cytostome, supported by microtubules. These are often packed together to form two or more rods, which function in ingestion, and in Entosiphon form an extendable siphon. Most phagotrophic euglenids have two flagella, one leading and one trailing. The latter is used for gliding along the substrate. In some, such as Peranema, the leading flagellum is rigid and beats only at its tip. Osmotrophic euglenoids Osmotrophic euglenids are euglenids which have undergone osmotrophy. Due to a lack of characteristics that are useful for taxonomical purposes, the origin of osmotrophic euglenids is unclear, though certain morphological characteristics reveal a small fraction of osmotrophic euglenids are derived from phototrophic and phagotrophic ancestors. A prolonged absence of light or exposure to harmful chemicals may cause atrophy and absorption of the chloroplasts without otherwise harming the organism. A number of species exists where a chloroplast's absence was formerly marked with separate genera such as Astasia (colourless Euglena) and Hyalophacus (colourless Phacus). Due to the lack of a developed cytostome, these forms feed exclusively by osmotrophic absorption. Reproduction Although euglenids share several common characteristics with animals, which is why they were originally classified as so, no evidence has been found of euglenids ever using sexual reproduction. This is one of the reasons they could no longer be classified as animals. For euglenids to reproduce, asexual reproduction takes place in the form of binary fission, and the cells replicate and divide during mitosis and cytokinesis. This process occurs in a very distinct order. First, the basal bodies and flagella replicate, then the cytostome and microtubules (the feeding apparatus), and finally the nucleus and remaining cytoskeleton. Once this occurs, the organism begins to cleave at the basal bodies, and this cleavage line moves towards the center of the organism until two separate euglenids are evident. Because of the way that this reproduction takes place and the axis of separation, it is called longitudinal cell division or longitudinal binary fission. Evolution The earliest fossil of euglenids is attributed to Moyeria, which is interpreted as possessing a pellicle composed of proteinaceous strips, the defining characteristic of euglenids. It is found in Middle Ordovician and Silurian rocks, making it the oldest fossil evidence of euglenids. Gallery	Biology and health sciences	Excavata	Plants
153008	https://en.wikipedia.org/wiki/Knot%20theory	Knot theory	In topology, knot theory is the study of mathematical knots. While inspired by knots which appear in daily life, such as those in shoelaces and rope, a mathematical knot differs in that the ends are joined so it cannot be undone, the simplest knot being a ring (or "unknot"). In mathematical language, a knot is an embedding of a circle in 3-dimensional Euclidean space, . Two mathematical knots are equivalent if one can be transformed into the other via a deformation of upon itself (known as an ambient isotopy); these transformations correspond to manipulations of a knotted string that do not involve cutting it or passing it through itself. Knots can be described in various ways. Using different description methods, there may be more than one description of the same knot. For example, a common method of describing a knot is a planar diagram called a knot diagram, in which any knot can be drawn in many different ways. Therefore, a fundamental problem in knot theory is determining when two descriptions represent the same knot. A complete algorithmic solution to this problem exists, which has unknown complexity. In practice, knots are often distinguished using a knot invariant, a "quantity" which is the same when computed from different descriptions of a knot. Important invariants include knot polynomials, knot groups, and hyperbolic invariants. The original motivation for the founders of knot theory was to create a table of knots and links, which are knots of several components entangled with each other. More than six billion knots and links have been tabulated since the beginnings of knot theory in the 19th century. To gain further insight, mathematicians have generalized the knot concept in several ways. Knots can be considered in other three-dimensional spaces and objects other than circles can be used; see knot (mathematics). For example, a higher-dimensional knot is an n-dimensional sphere embedded in (n+2)-dimensional Euclidean space. History Archaeologists have discovered that knot tying dates back to prehistoric times. Besides their uses such as recording information and tying objects together, knots have interested humans for their aesthetics and spiritual symbolism. Knots appear in various forms of Chinese artwork dating from several centuries BC (see Chinese knotting). The endless knot appears in Tibetan Buddhism, while the Borromean rings have made repeated appearances in different cultures, often representing strength in unity. The Celtic monks who created the Book of Kells lavished entire pages with intricate Celtic knotwork. A mathematical theory of knots was first developed in 1771 by Alexandre-Théophile Vandermonde who explicitly noted the importance of topological features when discussing the properties of knots related to the geometry of position. Mathematical studies of knots began in the 19th century with Carl Friedrich Gauss, who defined the linking integral . In the 1860s, Lord Kelvin's theory that atoms were knots in the aether led to Peter Guthrie Tait's creation of the first knot tables for complete classification. Tait, in 1885, published a table of knots with up to ten crossings, and what came to be known as the Tait conjectures. This record motivated the early knot theorists, but knot theory eventually became part of the emerging subject of topology. These topologists in the early part of the 20th century—Max Dehn, J. W. Alexander, and others—studied knots from the point of view of the knot group and invariants from homology theory such as the Alexander polynomial. This would be the main approach to knot theory until a series of breakthroughs transformed the subject. In the late 1970s, William Thurston introduced hyperbolic geometry into the study of knots with the hyperbolization theorem. Many knots were shown to be hyperbolic knots, enabling the use of geometry in defining new, powerful knot invariants. The discovery of the Jones polynomial by Vaughan Jones in 1984 , and subsequent contributions from Edward Witten, Maxim Kontsevich, and others, revealed deep connections between knot theory and mathematical methods in statistical mechanics and quantum field theory. A plethora of knot invariants have been invented since then, utilizing sophisticated tools such as quantum groups and Floer homology. In the last several decades of the 20th century, scientists became interested in studying physical knots in order to understand knotting phenomena in DNA and other polymers. Knot theory can be used to determine if a molecule is chiral (has a "handedness") or not . Tangles, strings with both ends fixed in place, have been effectively used in studying the action of topoisomerase on DNA . Knot theory may be crucial in the construction of quantum computers, through the model of topological quantum computation . Knot equivalence A knot is created by beginning with a one-dimensional line segment, wrapping it around itself arbitrarily, and then fusing its two free ends together to form a closed loop . Simply, we can say a knot is a "simple closed curve" (see Curve) — that is: a "nearly" injective and continuous function , with the only "non-injectivity" being . Topologists consider knots and other entanglements such as links and braids to be equivalent if the knot can be pushed about smoothly, without intersecting itself, to coincide with another knot. The idea of knot equivalence is to give a precise definition of when two knots should be considered the same even when positioned quite differently in space. A formal mathematical definition is that two knots are equivalent if there is an orientation-preserving homeomorphism with . What this definition of knot equivalence means is that two knots are equivalent when there is a continuous family of homeomorphisms of space onto itself, such that the last one of them carries the first knot onto the second knot. (In detail: Two knots and are equivalent if there exists a continuous mapping such that a) for each the mapping taking to is a homeomorphism of onto itself; b) for all ; and c) . Such a function is known as an ambient isotopy.) These two notions of knot equivalence agree exactly about which knots are equivalent: Two knots that are equivalent under the orientation-preserving homeomorphism definition are also equivalent under the ambient isotopy definition, because any orientation-preserving homeomorphisms of to itself is the final stage of an ambient isotopy starting from the identity. Conversely, two knots equivalent under the ambient isotopy definition are also equivalent under the orientation-preserving homeomorphism definition, because the (final) stage of the ambient isotopy must be an orientation-preserving homeomorphism carrying one knot to the other. The basic problem of knot theory, the recognition problem, is determining the equivalence of two knots. Algorithms exist to solve this problem, with the first given by Wolfgang Haken in the late 1960s . Nonetheless, these algorithms can be extremely time-consuming, and a major issue in the theory is to understand how hard this problem really is . The special case of recognizing the unknot, called the unknotting problem, is of particular interest . In February 2021 Marc Lackenby announced a new unknot recognition algorithm that runs in quasi-polynomial time. Knot diagrams A useful way to visualise and manipulate knots is to project the knot onto a plane—think of the knot casting a shadow on the wall. A small change in the direction of projection will ensure that it is one-to-one except at the double points, called crossings, where the "shadow" of the knot crosses itself once transversely . At each crossing, to be able to recreate the original knot, the over-strand must be distinguished from the under-strand. This is often done by creating a break in the strand going underneath. The resulting diagram is an immersed plane curve with the additional data of which strand is over and which is under at each crossing. (These diagrams are called knot diagrams when they represent a knot and link diagrams when they represent a link.) Analogously, knotted surfaces in 4-space can be related to immersed surfaces in 3-space. A reduced diagram is a knot diagram in which there are no reducible crossings (also nugatory or removable crossings), or in which all of the reducible crossings have been removed. A petal projection is a type of projection in which, instead of forming double points, all strands of the knot meet at a single crossing point, connected to it by loops forming non-nested "petals". Reidemeister moves In 1927, working with this diagrammatic form of knots, J. W. Alexander and Garland Baird Briggs, and independently Kurt Reidemeister, demonstrated that two knot diagrams belonging to the same knot can be related by a sequence of three kinds of moves on the diagram, shown below. These operations, now called the Reidemeister moves, are: The proof that diagrams of equivalent knots are connected by Reidemeister moves relies on an analysis of what happens under the planar projection of the movement taking one knot to another. The movement can be arranged so that almost all of the time the projection will be a knot diagram, except at finitely many times when an "event" or "catastrophe" occurs, such as when more than two strands cross at a point or multiple strands become tangent at a point. A close inspection will show that complicated events can be eliminated, leaving only the simplest events: (1) a "kink" forming or being straightened out; (2) two strands becoming tangent at a point and passing through; and (3) three strands crossing at a point. These are precisely the Reidemeister moves . Knot invariants A knot invariant is a "quantity" that is the same for equivalent knots . For example, if the invariant is computed from a knot diagram, it should give the same value for two knot diagrams representing equivalent knots. An invariant may take the same value on two different knots, so by itself may be incapable of distinguishing all knots. An elementary invariant is tricolorability. "Classical" knot invariants include the knot group, which is the fundamental group of the knot complement, and the Alexander polynomial, which can be computed from the Alexander invariant, a module constructed from the infinite cyclic cover of the knot complement . In the late 20th century, invariants such as "quantum" knot polynomials, Vassiliev invariants and hyperbolic invariants were discovered. These aforementioned invariants are only the tip of the iceberg of modern knot theory. Knot polynomials A knot polynomial is a knot invariant that is a polynomial. Well-known examples include the Jones polynomial, the Alexander polynomial, and the Kauffman polynomial. A variant of the Alexander polynomial, the Alexander–Conway polynomial, is a polynomial in the variable z with integer coefficients . The Alexander–Conway polynomial is actually defined in terms of links, which consist of one or more knots entangled with each other. The concepts explained above for knots, e.g. diagrams and Reidemeister moves, also hold for links. Consider an oriented link diagram, i.e. one in which every component of the link has a preferred direction indicated by an arrow. For a given crossing of the diagram, let be the oriented link diagrams resulting from changing the diagram as indicated in the figure: The original diagram might be either or , depending on the chosen crossing's configuration. Then the Alexander–Conway polynomial, , is recursively defined according to the rules: (where is any diagram of the unknot) The second rule is what is often referred to as a skein relation. To check that these rules give an invariant of an oriented link, one should determine that the polynomial does not change under the three Reidemeister moves. Many important knot polynomials can be defined in this way. The following is an example of a typical computation using a skein relation. It computes the Alexander–Conway polynomial of the trefoil knot. The yellow patches indicate where the relation is applied. C() = C() + z C() gives the unknot and the Hopf link. Applying the relation to the Hopf link where indicated, C() = C() + z C() gives a link deformable to one with 0 crossings (it is actually the unlink of two components) and an unknot. The unlink takes a bit of sneakiness: C() = C() + z C() which implies that C(unlink of two components) = 0, since the first two polynomials are of the unknot and thus equal. Putting all this together will show: Since the Alexander–Conway polynomial is a knot invariant, this shows that the trefoil is not equivalent to the unknot. So the trefoil really is "knotted". Actually, there are two trefoil knots, called the right and left-handed trefoils, which are mirror images of each other (take a diagram of the trefoil given above and change each crossing to the other way to get the mirror image). These are not equivalent to each other, meaning that they are not amphichiral. This was shown by Max Dehn, before the invention of knot polynomials, using group theoretical methods . But the Alexander–Conway polynomial of each kind of trefoil will be the same, as can be seen by going through the computation above with the mirror image. The Jones polynomial can in fact distinguish between the left- and right-handed trefoil knots . Hyperbolic invariants William Thurston proved many knots are hyperbolic knots, meaning that the knot complement (i.e., the set of points of 3-space not on the knot) admits a geometric structure, in particular that of hyperbolic geometry. The hyperbolic structure depends only on the knot so any quantity computed from the hyperbolic structure is then a knot invariant . Geometry lets us visualize what the inside of a knot or link complement looks like by imagining light rays as traveling along the geodesics of the geometry. An example is provided by the picture of the complement of the Borromean rings. The inhabitant of this link complement is viewing the space from near the red component. The balls in the picture are views of horoball neighborhoods of the link. By thickening the link in a standard way, the horoball neighborhoods of the link components are obtained. Even though the boundary of a neighborhood is a torus, when viewed from inside the link complement, it looks like a sphere. Each link component shows up as infinitely many spheres (of one color) as there are infinitely many light rays from the observer to the link component. The fundamental parallelogram (which is indicated in the picture), tiles both vertically and horizontally and shows how to extend the pattern of spheres infinitely. This pattern, the horoball pattern, is itself a useful invariant. Other hyperbolic invariants include the shape of the fundamental parallelogram, length of shortest geodesic, and volume. Modern knot and link tabulation efforts have utilized these invariants effectively. Fast computers and clever methods of obtaining these invariants make calculating these invariants, in practice, a simple task . Higher dimensions A knot in three dimensions can be untied when placed in four-dimensional space. This is done by changing crossings. Suppose one strand is behind another as seen from a chosen point. Lift it into the fourth dimension, so there is no obstacle (the front strand having no component there); then slide it forward, and drop it back, now in front. Analogies for the plane would be lifting a string up off the surface, or removing a dot from inside a circle. In fact, in four dimensions, any non-intersecting closed loop of one-dimensional string is equivalent to an unknot. First "push" the loop into a three-dimensional subspace, which is always possible, though technical to explain. Four-dimensional space occurs in classical knot theory, however, and an important topic is the study of slice knots and ribbon knots. A notorious open problem asks whether every slice knot is also ribbon. Knotting spheres of higher dimension Since a knot can be considered topologically a 1-dimensional sphere, the next generalization is to consider a two-dimensional sphere () embedded in 4-dimensional Euclidean space (). Such an embedding is knotted if there is no homeomorphism of onto itself taking the embedded 2-sphere to the standard "round" embedding of the 2-sphere. Suspended knots and spun knots are two typical families of such 2-sphere knots. The mathematical technique called "general position" implies that for a given n-sphere in m-dimensional Euclidean space, if m is large enough (depending on n), the sphere should be unknotted. In general, piecewise-linear n-spheres form knots only in (n + 2)-dimensional space , although this is no longer a requirement for smoothly knotted spheres. In fact, there are smoothly knotted -spheres in 6k-dimensional space; e.g., there is a smoothly knotted 3-sphere in . Thus the codimension of a smooth knot can be arbitrarily large when not fixing the dimension of the knotted sphere; however, any smooth k-sphere embedded in with is unknotted. The notion of a knot has further generalisations in mathematics, see: Knot (mathematics), isotopy classification of embeddings. Every knot in the n-sphere is the link of a real-algebraic set with isolated singularity in . An n-knot is a single embedded in . An n-link consists of k-copies of embedded in , where k is a natural number. Both the and the cases are well studied, and so is the case. Adding knots Two knots can be added by cutting both knots and joining the pairs of ends. The operation is called the knot sum, or sometimes the connected sum or composition of two knots. This can be formally defined as follows : consider a planar projection of each knot and suppose these projections are disjoint. Find a rectangle in the plane where one pair of opposite sides are arcs along each knot while the rest of the rectangle is disjoint from the knots. Form a new knot by deleting the first pair of opposite sides and adjoining the other pair of opposite sides. The resulting knot is a sum of the original knots. Depending on how this is done, two different knots (but no more) may result. This ambiguity in the sum can be eliminated regarding the knots as oriented, i.e. having a preferred direction of travel along the knot, and requiring the arcs of the knots in the sum are oriented consistently with the oriented boundary of the rectangle. The knot sum of oriented knots is commutative and associative. A knot is prime if it is non-trivial and cannot be written as the knot sum of two non-trivial knots. A knot that can be written as such a sum is composite. There is a prime decomposition for knots, analogous to prime and composite numbers . For oriented knots, this decomposition is also unique. Higher-dimensional knots can also be added but there are some differences. While you cannot form the unknot in three dimensions by adding two non-trivial knots, you can in higher dimensions, at least when one considers smooth knots in codimension at least 3. Knots can also be constructed using the circuit topology approach. This is done by combining basic units called soft contacts using five operations (Parallel, Series, Cross, Concerted, and Sub). The approach is applicable to open chains as well and can also be extended to include the so-called hard contacts. Tabulating knots Traditionally, knots have been catalogued in terms of crossing number. Knot tables generally include only prime knots, and only one entry for a knot and its mirror image (even if they are different) . The number of nontrivial knots of a given crossing number increases rapidly, making tabulation computationally difficult . Tabulation efforts have succeeded in enumerating over 6 billion knots and links . The sequence of the number of prime knots of a given crossing number, up to crossing number 16, is 0, 0, 1, 1, 2, 3, 7, 21, 49, 165, 552, 2176, 9988, , , ... . While exponential upper and lower bounds for this sequence are known, it has not been proven that this sequence is strictly increasing . The first knot tables by Tait, Little, and Kirkman used knot diagrams, although Tait also used a precursor to the Dowker notation. Different notations have been invented for knots which allow more efficient tabulation . The early tables attempted to list all knots of at most 10 crossings, and all alternating knots of 11 crossings . The development of knot theory due to Alexander, Reidemeister, Seifert, and others eased the task of verification and tables of knots up to and including 9 crossings were published by Alexander–Briggs and Reidemeister in the late 1920s. The first major verification of this work was done in the 1960s by John Horton Conway, who not only developed a new notation but also the Alexander–Conway polynomial . This verified the list of knots of at most 11 crossings and a new list of links up to 10 crossings. Conway found a number of omissions but only one duplication in the Tait–Little tables; however he missed the duplicates called the Perko pair, which would only be noticed in 1974 by Kenneth Perko . This famous error would propagate when Dale Rolfsen added a knot table in his influential text, based on Conway's work. Conway's 1970 paper on knot theory also contains a typographical duplication on its non-alternating 11-crossing knots page and omits 4 examples — 2 previously listed in D. Lombardero's 1968 Princeton senior thesis and 2 more subsequently discovered by Alain Caudron. [see Perko (1982), Primality of certain knots, Topology Proceedings] Less famous is the duplicate in his 10 crossing link table: 2.-2.-20.20 is the mirror of 8-20:-20. [See Perko (2016), Historical highlights of non-cyclic knot theory, J. Knot Theory Ramifications]. In the late 1990s Hoste, Thistlethwaite, and Weeks tabulated all the knots through 16 crossings . In 2003 Rankin, Flint, and Schermann, tabulated the alternating knots through 22 crossings . In 2020 Burton tabulated all prime knots with up to 19 crossings . Alexander–Briggs notation This is the most traditional notation, due to the 1927 paper of James W. Alexander and Garland B. Briggs and later extended by Dale Rolfsen in his knot table (see image above and List of prime knots). The notation simply organizes knots by their crossing number. One writes the crossing number with a subscript to denote its order amongst all knots with that crossing number. This order is arbitrary and so has no special significance (though in each number of crossings the twist knot comes after the torus knot). Links are written by the crossing number with a superscript to denote the number of components and a subscript to denote its order within the links with the same number of components and crossings. Thus the trefoil knot is notated 31 and the Hopf link is 2. Alexander–Briggs names in the range 10162 to 10166 are ambiguous, due to the discovery of the Perko pair in Charles Newton Little's original and subsequent knot tables, and differences in approach to correcting this error in knot tables and other publications created after this point. Dowker–Thistlethwaite notation The Dowker–Thistlethwaite notation, also called the Dowker notation or code, for a knot is a finite sequence of even integers. The numbers are generated by following the knot and marking the crossings with consecutive integers. Since each crossing is visited twice, this creates a pairing of even integers with odd integers. An appropriate sign is given to indicate over and undercrossing. For example, in this figure the knot diagram has crossings labelled with the pairs (1,6) (3,−12) (5,2) (7,8) (9,−4) and (11,−10). The Dowker–Thistlethwaite notation for this labelling is the sequence: 6, −12, 2, 8, −4, −10. A knot diagram has more than one possible Dowker notation, and there is a well-understood ambiguity when reconstructing a knot from a Dowker–Thistlethwaite notation. Conway notation The Conway notation for knots and links, named after John Horton Conway, is based on the theory of tangles . The advantage of this notation is that it reflects some properties of the knot or link. The notation describes how to construct a particular link diagram of the link. Start with a basic polyhedron, a 4-valent connected planar graph with no digon regions. Such a polyhedron is denoted first by the number of vertices then a number of asterisks which determine the polyhedron's position on a list of basic polyhedra. For example, 10* denotes the second 10-vertex polyhedron on Conway's list. Each vertex then has an algebraic tangle substituted into it (each vertex is oriented so there is no arbitrary choice in substitution). Each such tangle has a notation consisting of numbers and + or − signs. An example is 12 −3 2. The 1 denotes the only 1-vertex basic polyhedron. The 2 −3 2 is a sequence describing the continued fraction associated to a rational tangle. One inserts this tangle at the vertex of the basic polyhedron 1. A more complicated example is 83.1.2 0.1.1.1.1.1 Here again 8* refers to a basic polyhedron with 8 vertices. The periods separate the notation for each tangle. Any link admits such a description, and it is clear this is a very compact notation even for very large crossing number. There are some further shorthands usually used. The last example is usually written 83:2 0, where the ones are omitted and kept the number of dots excepting the dots at the end. For an algebraic knot such as in the first example, 1 is often omitted. Conway's pioneering paper on the subject lists up to 10-vertex basic polyhedra of which he uses to tabulate links, which have become standard for those links. For a further listing of higher vertex polyhedra, there are nonstandard choices available. Gauss code Gauss code, similar to the Dowker–Thistlethwaite notation, represents a knot with a sequence of integers. However, rather than every crossing being represented by two different numbers, crossings are labeled with only one number. When the crossing is an overcrossing, a positive number is listed. At an undercrossing, a negative number. For example, the trefoil knot in Gauss code can be given as: 1,−2,3,−1,2,−3 Gauss code is limited in its ability to identify knots. This problem is partially addressed with by the extended Gauss code.	Mathematics	Geometry	null
153221	https://en.wikipedia.org/wiki/Heat%20exchanger	Heat exchanger	A heat exchanger is a system used to transfer heat between a source and a working fluid. Heat exchangers are used in both cooling and heating processes. The fluids may be separated by a solid wall to prevent mixing or they may be in direct contact. They are widely used in space heating, refrigeration, air conditioning, power stations, chemical plants, petrochemical plants, petroleum refineries, natural-gas processing, and sewage treatment. The classic example of a heat exchanger is found in an internal combustion engine in which a circulating fluid known as engine coolant flows through radiator coils and air flows past the coils, which cools the coolant and heats the incoming air. Another example is the heat sink, which is a passive heat exchanger that transfers the heat generated by an electronic or a mechanical device to a fluid medium, often air or a liquid coolant. Flow arrangement There are three primary classifications of heat exchangers according to their flow arrangement. In parallel-flow heat exchangers, the two fluids enter the exchanger at the same end, and travel in parallel to one another to the other side. In counter-flow heat exchangers the fluids enter the exchanger from opposite ends. The counter current design is the most efficient, in that it can transfer the most heat from the heat (transfer) medium per unit mass due to the fact that the average temperature difference along any unit length is higher. See countercurrent exchange. In a cross-flow heat exchanger, the fluids travel roughly perpendicular to one another through the exchanger. For efficiency, heat exchangers are designed to maximize the surface area of the wall between the two fluids, while minimizing resistance to fluid flow through the exchanger. The exchanger's performance can also be affected by the addition of fins or corrugations in one or both directions, which increase surface area and may channel fluid flow or induce turbulence. The driving temperature across the heat transfer surface varies with position, but an appropriate mean temperature can be defined. In most simple systems this is the "log mean temperature difference" (LMTD). Sometimes direct knowledge of the LMTD is not available and the NTU method is used. Types By maximum operating temperature, heat exchangers can be divided into low-temperature and high-temperature ones. The former work up to 500–650°C depending on the industry and generally don't require special design and material considerations. The latter work up to 1000 or even 1400°C. Double pipe heat exchangers are the simplest exchangers used in industries. On one hand, these heat exchangers are cheap for both design and maintenance, making them a good choice for small industries. On the other hand, their low efficiency coupled with the high space occupied in large scales, has led modern industries to use more efficient heat exchangers like shell and tube or plate. However, since double pipe heat exchangers are simple, they are used to teach heat exchanger design basics to students as the fundamental rules for all heat exchangers are the same. 1. Double-pipe heat exchanger When one fluid flows through the smaller pipe, the other flows through the annular gap between the two pipes. These flows may be parallel or counter-flows in a double pipe heat exchanger. (a) Parallel flow, where both hot and cold liquids enter the heat exchanger from the same side, flow in the same direction and exit at the same end. This configuration is preferable when the two fluids are intended to reach exactly the same temperature, as it reduces thermal stress and produces a more uniform rate of heat transfer. (b) Counter-flow, where hot and cold fluids enter opposite sides of the heat exchanger, flow in opposite directions, and exit at opposite ends. This configuration is preferable when the objective is to maximize heat transfer between the fluids, as it creates a larger temperature differential when used under otherwise similar conditions. The figure above illustrates the parallel and counter-flow flow directions of the fluid exchanger. 2. Shell-and-tube heat exchanger In a shell-and-tube heat exchanger, two fluids at different temperatures flow through the heat exchanger. One of the fluids flows through the tube side and the other fluid flows outside the tubes, but inside the shell (shell side). Baffles are used to support the tubes, direct the fluid flow to the tubes in an approximately natural manner, and maximize the turbulence of the shell fluid. There are many various kinds of baffles, and the choice of baffle form, spacing, and geometry depends on the allowable flow rate of the drop in shell-side force, the need for tube support, and the flow-induced vibrations. There are several variations of shell-and-tube exchangers available; the differences lie in the arrangement of flow configurations and details of construction. In application to cool air with shell-and-tube technology (such as intercooler / charge air cooler for combustion engines), fins can be added on the tubes to increase heat transfer area on air side and create a tubes & fins configuration. 3. Plate Heat Exchanger A plate heat exchanger contains an amount of thin shaped heat transfer plates bundled together. The gasket arrangement of each pair of plates provides two separate channel system. Each pair of plates form a channel where the fluid can flow through. The pairs are attached by welding and bolting methods. The following shows the components in the heat exchanger. In single channels the configuration of the gaskets enables flow through. Thus, this allows the main and secondary media in counter-current flow. A gasket plate heat exchanger has a heat region from corrugated plates. The gasket function as seal between plates and they are located between frame and pressure plates. Fluid flows in a counter current direction throughout the heat exchanger. An efficient thermal performance is produced. Plates are produced in different depths, sizes and corrugated shapes. There are different types of plates available including plate and frame, plate and shell and spiral plate heat exchangers. The distribution area guarantees the flow of fluid to the whole heat transfer surface. This helps to prevent stagnant area that can cause accumulation of unwanted material on solid surfaces. High flow turbulence between plates results in a greater transfer of heat and a decrease in pressure. 4. Condensers and Boilers Heat exchangers using a two-phase heat transfer system are condensers, boilers and evaporators. Condensers are instruments that take and cool hot gas or vapor to the point of condensation and transform the gas into a liquid form. The point at which liquid transforms to gas is called vaporization and vice versa is called condensation. Surface condenser is the most common type of condenser where it includes a water supply device. Figure 5 below displays a two-pass surface condenser. The pressure of steam at the turbine outlet is low where the steam density is very low where the flow rate is very high. To prevent a decrease in pressure in the movement of steam from the turbine to condenser, the condenser unit is placed underneath and connected to the turbine. Inside the tubes the cooling water runs in a parallel way, while steam moves in a vertical downward position from the wide opening at the top and travel through the tube. Furthermore, boilers are categorized as initial application of heat exchangers. The word steam generator was regularly used to describe a boiler unit where a hot liquid stream is the source of heat rather than the combustion products. Depending on the dimensions and configurations the boilers are manufactured. Several boilers are only able to produce hot fluid while on the other hand the others are manufactured for steam production. Shell and tube Shell and tube heat exchangers consist of a series of tubes which contain fluid that must be either heated or cooled. A second fluid runs over the tubes that are being heated or cooled so that it can either provide the heat or absorb the heat required. A set of tubes is called the tube bundle and can be made up of several types of tubes: plain, longitudinally finned, etc. Shell and tube heat exchangers are typically used for high-pressure applications (with pressures greater than 30 bar and temperatures greater than 260 °C). This is because the shell and tube heat exchangers are robust due to their shape.Several thermal design features must be considered when designing the tubes in the shell and tube heat exchangers: There can be many variations on the shell and tube design. Typically, the ends of each tube are connected to plenums (sometimes called water boxes) through holes in tubesheets. The tubes may be straight or bent in the shape of a U, called U-tubes. Tube diameter: Using a small tube diameter makes the heat exchanger both economical and compact. However, it is more likely for the heat exchanger to foul up faster and the small size makes mechanical cleaning of the fouling difficult. To prevail over the fouling and cleaning problems, larger tube diameters can be used. Thus to determine the tube diameter, the available space, cost and fouling nature of the fluids must be considered. Tube thickness: The thickness of the wall of the tubes is usually determined to ensure: There is enough room for corrosion That flow-induced vibration has resistance Axial strength Availability of spare parts Hoop strength (to withstand internal tube pressure) Buckling strength (to withstand overpressure in the shell) Tube length: heat exchangers are usually cheaper when they have a smaller shell diameter and a long tube length. Thus, typically there is an aim to make the heat exchanger as long as physically possible whilst not exceeding production capabilities. However, there are many limitations for this, including space available at the installation site and the need to ensure tubes are available in lengths that are twice the required length (so they can be withdrawn and replaced). Also, long, thin tubes are difficult to take out and replace. Tube pitch: when designing the tubes, it is practical to ensure that the tube pitch (i.e., the centre-centre distance of adjoining tubes) is not less than 1.25 times the tubes' outside diameter. A larger tube pitch leads to a larger overall shell diameter, which leads to a more expensive heat exchanger. Tube corrugation: this type of tubes, mainly used for the inner tubes, increases the turbulence of the fluids and the effect is very important in the heat transfer giving a better performance. Tube Layout: refers to how tubes are positioned within the shell. There are four main types of tube layout, which are, triangular (30°), rotated triangular (60°), square (90°) and rotated square (45°). The triangular patterns are employed to give greater heat transfer as they force the fluid to flow in a more turbulent fashion around the piping. Square patterns are employed where high fouling is experienced and cleaning is more regular. Baffle Design: baffles are used in shell and tube heat exchangers to direct fluid across the tube bundle. They run perpendicularly to the shell and hold the bundle, preventing the tubes from sagging over a long length. They can also prevent the tubes from vibrating. The most common type of baffle is the segmental baffle. The semicircular segmental baffles are oriented at 180 degrees to the adjacent baffles forcing the fluid to flow upward and downwards between the tube bundle. Baffle spacing is of large thermodynamic concern when designing shell and tube heat exchangers. Baffles must be spaced with consideration for the conversion of pressure drop and heat transfer. For thermo economic optimization it is suggested that the baffles be spaced no closer than 20% of the shell's inner diameter. Having baffles spaced too closely causes a greater pressure drop because of flow redirection. Consequently, having the baffles spaced too far apart means that there may be cooler spots in the corners between baffles. It is also important to ensure the baffles are spaced close enough that the tubes do not sag. The other main type of baffle is the disc and doughnut baffle, which consists of two concentric baffles. An outer, wider baffle looks like a doughnut, whilst the inner baffle is shaped like a disk. This type of baffle forces the fluid to pass around each side of the disk then through the doughnut baffle generating a different type of fluid flow. Tubes & fins Design: in application to cool air with shell-and-tube technology (such as intercooler / charge air cooler for combustion engines), the difference in heat transfer between air and cold fluid can be such that there is a need to increase heat transfer area on air side. For this function fins can be added on the tubes to increase heat transfer area on air side and create a tubes & fins configuration. Fixed tube liquid-cooled heat exchangers especially suitable for marine and harsh applications can be assembled with brass shells, copper tubes, brass baffles, and forged brass integral end hubs. (See: Copper in heat exchangers). Plate Another type of heat exchanger is the plate heat exchanger. These exchangers are composed of many thin, slightly separated plates that have very large surface areas and small fluid flow passages for heat transfer. Advances in gasket and brazing technology have made the plate-type heat exchanger increasingly practical. In HVAC applications, large heat exchangers of this type are called plate-and-frame; when used in open loops, these heat exchangers are normally of the gasket type to allow periodic disassembly, cleaning, and inspection. There are many types of permanently bonded plate heat exchangers, such as dip-brazed, vacuum-brazed, and welded plate varieties, and they are often specified for closed-loop applications such as refrigeration. Plate heat exchangers also differ in the types of plates that are used, and in the configurations of those plates. Some plates may be stamped with "chevron", dimpled, or other patterns, where others may have machined fins and/or grooves. When compared to shell and tube exchangers, the stacked-plate arrangement typically has lower volume and cost. Another difference between the two is that plate exchangers typically serve low to medium pressure fluids, compared to medium and high pressures of shell and tube. A third and important difference is that plate exchangers employ more countercurrent flow rather than cross current flow, which allows lower approach temperature differences, high temperature changes, and increased efficiencies. Plate and shell A third type of heat exchanger is a plate and shell heat exchanger, which combines plate heat exchanger with shell and tube heat exchanger technologies. The heart of the heat exchanger contains a fully welded circular plate pack made by pressing and cutting round plates and welding them together. Nozzles carry flow in and out of the platepack (the 'Plate side' flowpath). The fully welded platepack is assembled into an outer shell that creates a second flowpath (the 'Shell side'). Plate and shell technology offers high heat transfer, high pressure, high operating temperature, compact size, low fouling and close approach temperature. In particular, it does completely without gaskets, which provides security against leakage at high pressures and temperatures. Adiabatic wheel A fourth type of heat exchanger uses an intermediate fluid or solid store to hold heat, which is then moved to the other side of the heat exchanger to be released. Two examples of this are adiabatic wheels, which consist of a large wheel with fine threads rotating through the hot and cold fluids, and fluid heat exchangers. Plate fin This type of heat exchanger uses "sandwiched" passages containing fins to increase the effectiveness of the unit. The designs include crossflow and counterflow coupled with various fin configurations such as straight fins, offset fins and wavy fins. Plate and fin heat exchangers are usually made of aluminum alloys, which provide high heat transfer efficiency. The material enables the system to operate at a lower temperature difference and reduce the weight of the equipment. Plate and fin heat exchangers are mostly used for low temperature services such as natural gas, helium and oxygen liquefaction plants, air separation plants and transport industries such as motor and aircraft engines. Advantages of plate and fin heat exchangers: High heat transfer efficiency especially in gas treatment Larger heat transfer area Approximately 5 times lighter in weight than that of shell and tube heat exchanger. Able to withstand high pressure Disadvantages of plate and fin heat exchangers: Might cause clogging as the pathways are very narrow Difficult to clean the pathways Aluminium alloys are susceptible to Mercury Liquid Embrittlement Failure Finned tube The usage of fins in a tube-based heat exchanger is common when one of the working fluids is a low-pressure gas, and is typical for heat exchangers that operate using ambient air, such as automotive radiators and HVAC air condensers. Fins dramatically increase the surface area with which heat can be exchanged, which improves the efficiency of conducting heat to a fluid with very low thermal conductivity, such as air. The fins are typically made from aluminium or copper since they must conduct heat from the tube along the length of the fins, which are usually very thin. The main construction types of finned tube exchangers are: A stack of evenly-spaced metal plates act as the fins and the tubes are pressed through pre-cut holes in the fins, good thermal contact usually being achieved by deformation of the fins around the tube. This is typical construction for HVAC air coils and large refrigeration condensers. Fins are spiral-wound onto individual tubes as a continuous strip, the tubes can then be assembled in banks, bent in a serpentine pattern, or wound into large spirals. Zig-zag metal strips are sandwiched between flat rectangular tubes, often being soldered or brazed together for good thermal and mechanical strength. This is common in low-pressure heat exchangers such as water-cooling radiators. Regular flat tubes will expand and deform if exposed to high pressures but flat microchannel tubes allow this construction to be used for high pressures. Stacked-fin or spiral-wound construction can be used for the tubes inside shell-and-tube heat exchangers when high efficiency thermal transfer to a gas is required. In electronics cooling, heat sinks, particularly those using heat pipes, can have a stacked-fin construction. Pillow plate A pillow plate heat exchanger is commonly used in the dairy industry for cooling milk in large direct-expansion stainless steel bulk tanks. Nearly the entire surface area of a tank can be integrated with this heat exchanger, without gaps that would occur between pipes welded to the exterior of the tank. Pillow plates can also be constructed as flat plates that are stacked inside a tank. The relatively flat surface of the plates allows easy cleaning, especially in sterile applications. The pillow plate can be constructed using either a thin sheet of metal welded to the thicker surface of a tank or vessel, or two thin sheets welded together. The surface of the plate is welded with a regular pattern of dots or a serpentine pattern of weld lines. After welding the enclosed space is pressurised with sufficient force to cause the thin metal to bulge out around the welds, providing a space for heat exchanger liquids to flow, and creating a characteristic appearance of a swelled pillow formed out of metal. Waste heat recovery units A waste heat recovery unit (WHRU) is a heat exchanger that recovers heat from a hot gas stream while transferring it to a working medium, typically water or oils. The hot gas stream can be the exhaust gas from a gas turbine or a diesel engine or a waste gas from industry or refinery. Large systems with high volume and temperature gas streams, typical in industry, can benefit from steam Rankine cycle (SRC) in a waste heat recovery unit, but these cycles are too expensive for small systems. The recovery of heat from low temperature systems requires different working fluids than steam. An organic Rankine cycle (ORC) waste heat recovery unit can be more efficient at low temperature range using refrigerants that boil at lower temperatures than water. Typical organic refrigerants are ammonia, pentafluoropropane (R-245fa and R-245ca), and toluene. The refrigerant is boiled by the heat source in the evaporator to produce super-heated vapor. This fluid is expanded in the turbine to convert thermal energy to kinetic energy, that is converted to electricity in the electrical generator. This energy transfer process decreases the temperature of the refrigerant that, in turn, condenses. The cycle is closed and completed using a pump to send the fluid back to the evaporator. Dynamic scraped surface Another type of heat exchanger is called "(dynamic) scraped surface heat exchanger". This is mainly used for heating or cooling with high-viscosity products, crystallization processes, evaporation and high-fouling applications. Long running times are achieved due to the continuous scraping of the surface, thus avoiding fouling and achieving a sustainable heat transfer rate during the process. Phase-change In addition to heating up or cooling down fluids in just a single phase, heat exchangers can be used either to heat a liquid to evaporate (or boil) it or used as condensers to cool a vapor and condense it to a liquid. In chemical plants and refineries, reboilers used to heat incoming feed for distillation towers are often heat exchangers. Distillation set-ups typically use condensers to condense distillate vapors back into liquid. Power plants that use steam-driven turbines commonly use heat exchangers to boil water into steam. Heat exchangers or similar units for producing steam from water are often called boilers or steam generators. In the nuclear power plants called pressurized water reactors, special large heat exchangers pass heat from the primary (reactor plant) system to the secondary (steam plant) system, producing steam from water in the process. These are called steam generators. All fossil-fueled and nuclear power plants using steam-driven turbines have surface condensers to convert the exhaust steam from the turbines into condensate (water) for re-use. To conserve energy and cooling capacity in chemical and other plants, regenerative heat exchangers can transfer heat from a stream that must be cooled to another stream that must be heated, such as distillate cooling and reboiler feed pre-heating. This term can also refer to heat exchangers that contain a material within their structure that has a change of phase. This is usually a solid to liquid phase due to the small volume difference between these states. This change of phase effectively acts as a buffer because it occurs at a constant temperature but still allows for the heat exchanger to accept additional heat. One example where this has been investigated is for use in high power aircraft electronics. Heat exchangers functioning in multiphase flow regimes may be subject to the Ledinegg instability. Direct contact Direct contact heat exchangers involve heat transfer between hot and cold streams of two phases in the absence of a separating wall. Thus such heat exchangers can be classified as: Gas – liquid Immiscible liquid – liquid Solid-liquid or solid – gas Most direct contact heat exchangers fall under the Gas – Liquid category, where heat is transferred between a gas and liquid in the form of drops, films or sprays. Such types of heat exchangers are used predominantly in air conditioning, humidification, industrial hot water heating, water cooling and condensing plants. Microchannel Microchannel heat exchangers are multi-pass parallel flow heat exchangers consisting of three main elements: manifolds (inlet and outlet), multi-port tubes with the hydraulic diameters smaller than 1mm, and fins. All the elements usually brazed together using controllable atmosphere brazing process. Microchannel heat exchangers are characterized by high heat transfer ratio, low refrigerant charges, compact size, and lower airside pressure drops compared to finned tube heat exchangers. Microchannel heat exchangers are widely used in automotive industry as the car radiators, and as condenser, evaporator, and cooling/heating coils in HVAC industry. Micro heat exchangers, Micro-scale heat exchangers, or microstructured heat exchangers are heat exchangers in which (at least one) fluid flows in lateral confinements with typical dimensions below 1 mm. The most typical such confinement are microchannels, which are channels with a hydraulic diameter below 1 mm. Microchannel heat exchangers can be made from metal or ceramics. Microchannel heat exchangers can be used for many applications including: high-performance aircraft gas turbine engines heat pumps Microprocessor and microchip cooling air conditioning HVAC and refrigeration air coils One of the widest uses of heat exchangers is for refrigeration and air conditioning. This class of heat exchangers is commonly called air coils, or just coils due to their often-serpentine internal tubing, or condensers in the case of refrigeration, and are typically of the finned tube type. Liquid-to-air, or air-to-liquid HVAC coils are typically of modified crossflow arrangement. In vehicles, heat coils are often called heater cores. On the liquid side of these heat exchangers, the common fluids are water, a water-glycol solution, steam, or a refrigerant. For heating coils, hot water and steam are the most common, and this heated fluid is supplied by boilers, for example. For cooling coils, chilled water and refrigerant are most common. Chilled water is supplied from a chiller that is potentially located very far away, but refrigerant must come from a nearby condensing unit. When a refrigerant is used, the cooling coil is the evaporator, and the heating coil is the condenser in the vapor-compression refrigeration cycle. HVAC coils that use this direct-expansion of refrigerants are commonly called DX coils. Some DX coils are "microchannel" type. On the air side of HVAC coils a significant difference exists between those used for heating, and those for cooling. Due to psychrometrics, air that is cooled often has moisture condensing out of it, except with extremely dry air flows. Heating some air increases that airflow's capacity to hold water. So heating coils need not consider moisture condensation on their air-side, but cooling coils must be adequately designed and selected to handle their particular latent (moisture) as well as the sensible (cooling) loads. The water that is removed is called condensate. For many climates, water or steam HVAC coils can be exposed to freezing conditions. Because water expands upon freezing, these somewhat expensive and difficult to replace thin-walled heat exchangers can easily be damaged or destroyed by just one freeze. As such, freeze protection of coils is a major concern of HVAC designers, installers, and operators. The introduction of indentations placed within the heat exchange fins controlled condensation, allowing water molecules to remain in the cooled air. The heat exchangers in direct-combustion furnaces, typical in many residences, are not 'coils'. They are, instead, gas-to-air heat exchangers that are typically made of stamped steel sheet metal. The combustion products pass on one side of these heat exchangers, and air to heat on the other. A cracked heat exchanger is therefore a dangerous situation that requires immediate attention because combustion products may enter living space. Helical-coil Although double-pipe heat exchangers are the simplest to design, the better choice in the following cases would be the helical-coil heat exchanger (HCHE): The main advantage of the HCHE, like that for the Spiral heat exchanger (SHE), is its highly efficient use of space, especially when it's limited and not enough straight pipe can be laid. Under conditions of low flowrates (or laminar flow), such that the typical shell-and-tube exchangers have low heat-transfer coefficients and becoming uneconomical. When there is low pressure in one of the fluids, usually from accumulated pressure drops in other process equipment. When one of the fluids has components in multiple phases (solids, liquids, and gases), which tends to create mechanical problems during operations, such as plugging of small-diameter tubes. Cleaning of helical coils for these multiple-phase fluids can prove to be more difficult than its shell and tube counterpart; however the helical coil unit would require cleaning less often. These have been used in the nuclear industry as a method for exchanging heat in a sodium system for large liquid metal fast breeder reactors since the early 1970s, using an HCHE device invented by Charles E. Boardman and John H. Germer. There are several simple methods for designing HCHE for all types of manufacturing industries, such as using the Ramachandra K. Patil (et al.) method from India and the Scott S. Haraburda method from the United States. However, these are based upon assumptions of estimating inside heat transfer coefficient, predicting flow around the outside of the coil, and upon constant heat flux. Spiral A modification to the perpendicular flow of the typical HCHE involves the replacement of shell with another coiled tube, allowing the two fluids to flow parallel to one another, and which requires the use of different design calculations. These are the Spiral Heat Exchangers (SHE), which may refer to a helical (coiled) tube configuration, more generally, the term refers to a pair of flat surfaces that are coiled to form the two channels in a counter-flow arrangement. Each of the two channels has one long curved path. A pair of fluid ports are connected tangentially to the outer arms of the spiral, and axial ports are common, but optional. The main advantage of the SHE is its highly efficient use of space. This attribute is often leveraged and partially reallocated to gain other improvements in performance, according to well known tradeoffs in heat exchanger design. (A notable tradeoff is capital cost vs operating cost.) A compact SHE may be used to have a smaller footprint and thus lower all-around capital costs, or an oversized SHE may be used to have less pressure drop, less pumping energy, higher thermal efficiency, and lower energy costs. Construction The distance between the sheets in the spiral channels is maintained by using spacer studs that were welded prior to rolling. Once the main spiral pack has been rolled, alternate top and bottom edges are welded and each end closed by a gasketed flat or conical cover bolted to the body. This ensures no mixing of the two fluids occurs. Any leakage is from the periphery cover to the atmosphere, or to a passage that contains the same fluid. Self cleaning Spiral heat exchangers are often used in the heating of fluids that contain solids and thus tend to foul the inside of the heat exchanger. The low pressure drop lets the SHE handle fouling more easily. The SHE uses a “self cleaning” mechanism, whereby fouled surfaces cause a localized increase in fluid velocity, thus increasing the drag (or fluid friction) on the fouled surface, thus helping to dislodge the blockage and keep the heat exchanger clean. "The internal walls that make up the heat transfer surface are often rather thick, which makes the SHE very robust, and able to last a long time in demanding environments." They are also easily cleaned, opening out like an oven where any buildup of foulant can be removed by pressure washing. Self-cleaning water filters are used to keep the system clean and running without the need to shut down or replace cartridges and bags. Flow arrangements There are three main types of flows in a spiral heat exchanger: Counter-current Flow: Fluids flow in opposite directions. These are used for liquid-liquid, condensing and gas cooling applications. Units are usually mounted vertically when condensing vapour and mounted horizontally when handling high concentrations of solids. Spiral Flow/Cross Flow: One fluid is in spiral flow and the other in a cross flow. Spiral flow passages are welded at each side for this type of spiral heat exchanger. This type of flow is suitable for handling low density gas, which passes through the cross flow, avoiding pressure loss. It can be used for liquid-liquid applications if one liquid has a considerably greater flow rate than the other. Distributed Vapour/Spiral flow: This design is that of a condenser, and is usually mounted vertically. It is designed to cater for the sub-cooling of both condensate and non-condensables. The coolant moves in a spiral and leaves via the top. Hot gases that enter leave as condensate via the bottom outlet. Applications The Spiral heat exchanger is good for applications such as pasteurization, digester heating, heat recovery, pre-heating (see: recuperator), and effluent cooling. For sludge treatment, SHEs are generally smaller than other types of heat exchangers. These are used to transfer the heat. Selection Due to the many variables involved, selecting optimal heat exchangers is challenging. Hand calculations are possible, but many iterations are typically needed. As such, heat exchangers are most often selected via computer programs, either by system designers, who are typically engineers, or by equipment vendors. To select an appropriate heat exchanger, the system designers (or equipment vendors) would firstly consider the design limitations for each heat exchanger type. Though cost is often the primary criterion, several other selection criteria are important: High/low pressure limits Thermal performance Temperature ranges Product mix (liquid/liquid, particulates or high-solids liquid) Pressure drops across the exchanger Fluid flow capacity Cleanability, maintenance and repair Materials required for construction Ability and ease of future expansion Material selection, such as copper, aluminium, carbon steel, stainless steel, nickel alloys, ceramic, polymer, and titanium. Small-diameter coil technologies are becoming more popular in modern air conditioning and refrigeration systems because they have better rates of heat transfer than conventional sized condenser and evaporator coils with round copper tubes and aluminum or copper fin that have been the standard in the HVAC industry. Small diameter coils can withstand the higher pressures required by the new generation of environmentally friendlier refrigerants. Two small diameter coil technologies are currently available for air conditioning and refrigeration products: copper microgroove and brazed aluminum microchannel. Choosing the right heat exchanger (HX) requires some knowledge of the different heat exchanger types, as well as the environment where the unit must operate. Typically in the manufacturing industry, several differing types of heat exchangers are used for just one process or system to derive the final product. For example, a kettle HX for pre-heating, a double pipe HX for the 'carrier' fluid and a plate and frame HX for final cooling. With sufficient knowledge of heat exchanger types and operating requirements, an appropriate selection can be made to optimise the process. Monitoring and maintenance Online monitoring of commercial heat exchangers is done by tracking the overall heat transfer coefficient. The overall heat transfer coefficient tends to decline over time due to fouling. By periodically calculating the overall heat transfer coefficient from exchanger flow rates and temperatures, the owner of the heat exchanger can estimate when cleaning the heat exchanger is economically attractive. Integrity inspection of plate and tubular heat exchanger can be tested in situ by the conductivity or helium gas methods. These methods confirm the integrity of the plates or tubes to prevent any cross contamination and the condition of the gaskets. Mechanical integrity monitoring of heat exchanger tubes may be conducted through Nondestructive methods such as eddy current testing. Fouling Fouling occurs when impurities deposit on the heat exchange surface. Deposition of these impurities can decrease heat transfer effectiveness significantly over time and are caused by: Low wall shear stress Low fluid velocities High fluid velocities Reaction product solid precipitation Precipitation of dissolved impurities due to elevated wall temperatures The rate of heat exchanger fouling is determined by the rate of particle deposition less re-entrainment/suppression. This model was originally proposed in 1959 by Kern and Seaton. Crude Oil Exchanger Fouling. In commercial crude oil refining, crude oil is heated from to prior to entering the distillation column. A series of shell and tube heat exchangers typically exchange heat between crude oil and other oil streams to heat the crude to prior to heating in a furnace. Fouling occurs on the crude side of these exchangers due to asphaltene insolubility. The nature of asphaltene solubility in crude oil was successfully modeled by Wiehe and Kennedy. The precipitation of insoluble asphaltenes in crude preheat trains has been successfully modeled as a first order reaction by Ebert and Panchal who expanded on the work of Kern and Seaton. Cooling Water Fouling. Cooling water systems are susceptible to fouling. Cooling water typically has a high total dissolved solids content and suspended colloidal solids. Localized precipitation of dissolved solids occurs at the heat exchange surface due to wall temperatures higher than bulk fluid temperature. Low fluid velocities (less than 3 ft/s) allow suspended solids to settle on the heat exchange surface. Cooling water is typically on the tube side of a shell and tube exchanger because it's easy to clean. To prevent fouling, designers typically ensure that cooling water velocity is greater than and bulk fluid temperature is maintained less than . Other approaches to control fouling control combine the "blind" application of biocides and anti-scale chemicals with periodic lab testing. Maintenance Plate and frame heat exchangers can be disassembled and cleaned periodically. Tubular heat exchangers can be cleaned by such methods as acid cleaning, sandblasting, high-pressure water jet, bullet cleaning, or drill rods. In large-scale cooling water systems for heat exchangers, water treatment such as purification, addition of chemicals, and testing, is used to minimize fouling of the heat exchange equipment. Other water treatment is also used in steam systems for power plants, etc. to minimize fouling and corrosion of the heat exchange and other equipment. A variety of companies have started using water borne oscillations technology to prevent biofouling. Without the use of chemicals, this type of technology has helped in providing a low-pressure drop in heat exchangers. Design and manufacturing regulations The design and manufacturing of heat exchangers has numerous regulations, which vary according to the region in which they will be used. Design and manufacturing codes include: ASME Boiler and Pressure Vessel Code (US); PD 5500 (UK); BS 1566 (UK); EN 13445 (EU); CODAP (French); Pressure Equipment Safety Regulations 2016 (PER) (UK); Pressure Equipment Directive (EU); NORSOK (Norwegian); TEMA; API 12; and API 560. In nature Humans The human nasal passages serve as a heat exchanger, with cool air being inhaled and warm air being exhaled. Its effectiveness can be demonstrated by putting the hand in front of the face and exhaling, first through the nose and then through the mouth. Air exhaled through the nose is substantially cooler. This effect can be enhanced with clothing, by, for example, wearing a scarf over the face while breathing in cold weather. In species that have external testes (such as human), the artery to the testis is surrounded by a mesh of veins called the pampiniform plexus. This cools the blood heading to the testes, while reheating the returning blood. Birds, fish, marine mammals "Countercurrent" heat exchangers occur naturally in the circulatory systems of fish, whales and other marine mammals. Arteries to the skin carrying warm blood are intertwined with veins from the skin carrying cold blood, causing the warm arterial blood to exchange heat with the cold venous blood. This reduces the overall heat loss in cold water. Heat exchangers are also present in the tongues of baleen whales as large volumes of water flow through their mouths. Wading birds use a similar system to limit heat losses from their body through their legs into the water. Carotid rete Carotid rete is a counter-current heat exchanging organ in some ungulates. The blood ascending the carotid arteries on its way to the brain, flows via a network of vessels where heat is discharged to the veins of cooler blood descending from the nasal passages. The carotid rete allows Thomson's gazelle to maintain its brain almost 3 °C (5.4 °F) cooler than the rest of the body, and therefore aids in tolerating bursts in metabolic heat production such as associated with outrunning cheetahs (during which the body temperature exceeds the maximum temperature at which the brain could function). Humans with other primates lack a carotid rete. In industry Heat exchangers are widely used in industry both for cooling and heating large scale industrial processes. The type and size of heat exchanger used can be tailored to suit a process depending on the type of fluid, its phase, temperature, density, viscosity, pressures, chemical composition and various other thermodynamic properties. In many industrial processes there is waste of energy or a heat stream that is being exhausted, heat exchangers can be used to recover this heat and put it to use by heating a different stream in the process. This practice saves a lot of money in industry, as the heat supplied to other streams from the heat exchangers would otherwise come from an external source that is more expensive and more harmful to the environment. Heat exchangers are used in many industries, including: Waste water treatment Refrigeration Wine and beer making Petroleum refining Nuclear power In waste water treatment, heat exchangers play a vital role in maintaining optimal temperatures within anaerobic digesters to promote the growth of microbes that remove pollutants. Common types of heat exchangers used in this application are the double pipe heat exchanger as well as the plate and frame heat exchanger. In aircraft In commercial aircraft heat exchangers are used to take heat from the engine's oil system to heat cold fuel. This improves fuel efficiency, as well as reduces the possibility of water entrapped in the fuel freezing in components. Current market and forecast Estimated at US$17.5 billion in 2021, the global demand of heat exchangers is expected to experience robust growth of about 5% annually over the next years. The market value is expected to reach US$27 billion by 2030. With an expanding desire for environmentally friendly options and increased development of offices, retail sectors, and public buildings, market expansion is due to grow. A model of a simple heat exchanger A simple heat exchange might be thought of as two straight pipes with fluid flow, which are thermally connected. Let the pipes be of equal length L, carrying fluids with heat capacity (energy per unit mass per unit change in temperature) and let the mass flow rate of the fluids through the pipes, both in the same direction, be (mass per unit time), where the subscript i applies to pipe 1 or pipe 2. Temperature profiles for the pipes are and where x is the distance along the pipe. Assume a steady state, so that the temperature profiles are not functions of time. Assume also that the only transfer of heat from a small volume of fluid in one pipe is to the fluid element in the other pipe at the same position, i.e., there is no transfer of heat along a pipe due to temperature differences in that pipe. By Newton's law of cooling the rate of change in energy of a small volume of fluid is proportional to the difference in temperatures between it and the corresponding element in the other pipe: ( this is for parallel flow in the same direction and opposite temperature gradients, but for counter-flow heat exchange countercurrent exchange the sign is opposite in the second equation in front of ), where is the thermal energy per unit length and γ is the thermal connection constant per unit length between the two pipes. This change in internal energy results in a change in the temperature of the fluid element. The time rate of change for the fluid element being carried along by the flow is: where is the "thermal mass flow rate". The differential equations governing the heat exchanger may now be written as: Since the system is in a steady state, there are no partial derivatives of temperature with respect to time, and since there is no heat transfer along the pipe, there are no second derivatives in x as is found in the heat equation. These two coupled first-order differential equations may be solved to yield: where , , (this is for parallel-flow, but for counter-flow the sign in front of is negative, so that if , for the same "thermal mass flow rate" in both opposite directions, the gradient of temperature is constant and the temperatures linear in position x with a constant difference along the exchanger, explaining why the counter current design countercurrent exchange is the most efficient ) and A and B are two as yet undetermined constants of integration. Let and be the temperatures at x=0 and let and be the temperatures at the end of the pipe at x=L. Define the average temperatures in each pipe as: Using the solutions above, these temperatures are: {\| \|- \| \| \|- \| \| \|- \| \| \|} Choosing any two of the temperatures above eliminates the constants of integration, letting us find the other four temperatures. We find the total energy transferred by integrating the expressions for the time rate of change of internal energy per unit length: By the conservation of energy, the sum of the two energies is zero. The quantity is known as the Log mean temperature difference, and is a measure of the effectiveness of the heat exchanger in transferring heat energy.	Technology	Heating and cooling	null
153299	https://en.wikipedia.org/wiki/Root-finding%20algorithm	Root-finding algorithm	In numerical analysis, a root-finding algorithm is an algorithm for finding zeros, also called "roots", of continuous functions. A zero of a function is a number such that . As, generally, the zeros of a function cannot be computed exactly nor expressed in closed form, root-finding algorithms provide approximations to zeros. For functions from the real numbers to real numbers or from the complex numbers to the complex numbers, these are expressed either as floating-point numbers without error bounds or as floating-point values together with error bounds. The latter, approximations with error bounds, are equivalent to small isolating intervals for real roots or disks for complex roots. Solving an equation is the same as finding the roots of the function . Thus root-finding algorithms can be used to solve any equation of continuous functions. However, most root-finding algorithms do not guarantee that they will find all roots of a function, and if such an algorithm does not find any root, that does not necessarily mean that no root exists. Most numerical root-finding methods are iterative methods, producing a sequence of numbers that ideally converges towards a root as a limit. They require one or more initial guesses of the root as starting values, then each iteration of the algorithm produces a successively more accurate approximation to the root. Since the iteration must be stopped at some point, these methods produce an approximation to the root, not an exact solution. Many methods compute subsequent values by evaluating an auxiliary function on the preceding values. The limit is thus a fixed point of the auxiliary function, which is chosen for having the roots of the original equation as fixed points and for converging rapidly to these fixed points. The behavior of general root-finding algorithms is studied in numerical analysis. However, for polynomials specifically, the study of root-finding algorithms belongs to computer algebra, since algebraic properties of polynomials are fundamental for the most efficient algorithms. The efficiency and applicability of an algorithm may depend sensitively on the characteristics of the given functions. For example, many algorithms use the derivative of the input function, while others work on every continuous function. In general, numerical algorithms are not guaranteed to find all the roots of a function, so failing to find a root does not prove that there is no root. However, for polynomials, there are specific algorithms that use algebraic properties for certifying that no root is missed and for locating the roots in separate intervals (or disks for complex roots) that are small enough to ensure the convergence of numerical methods (typically Newton's method) to the unique root within each interval (or disk). Bracketing methods Bracketing methods determine successively smaller intervals (brackets) that contain a root. When the interval is small enough, then a root is considered found. These generally use the intermediate value theorem, which asserts that if a continuous function has values of opposite signs at the end points of an interval, then the function has at least one root in the interval. Therefore, they require starting with an interval such that the function takes opposite signs at the end points of the interval. However, in the case of polynomials there are other methods such as Descartes' rule of signs, Budan's theorem and Sturm's theorem for bounding or determining the number of roots in an interval. They lead to efficient algorithms for real-root isolation of polynomials, which find all real roots with a guaranteed accuracy. Bisection method The simplest root-finding algorithm is the bisection method. Let be a continuous function for which one knows an interval such that and have opposite signs (a bracket). Let be the middle of the interval (the midpoint or the point that bisects the interval). Then either and , or and have opposite signs, and one has divided by two the size of the interval. Although the bisection method is robust, it gains one and only one bit of accuracy with each iteration. Therefore, the number of function evaluations required for finding an ε-approximate root is . Other methods, under appropriate conditions, can gain accuracy faster. False position (regula falsi) The false position method, also called the regula falsi method, is similar to the bisection method, but instead of using bisection search's middle of the interval it uses the -intercept of the line that connects the plotted function values at the endpoints of the interval, that is False position is similar to the secant method, except that, instead of retaining the last two points, it makes sure to keep one point on either side of the root. The false position method can be faster than the bisection method and will never diverge like the secant method. However, it may fail to converge in some naive implementations due to roundoff errors that may lead to a wrong sign for . Typically, this may occur if the derivative of is large in the neighborhood of the root. ITP method The ITP method is the only known method to bracket the root with the same worst case guarantees of the bisection method while guaranteeing a superlinear convergence to the root of smooth functions as the secant method. It is also the only known method guaranteed to outperform the bisection method on the average for any continuous distribution on the location of the root (see ITP Method#Analysis). It does so by keeping track of both the bracketing interval as well as the minmax interval in which any point therein converges as fast as the bisection method. The construction of the queried point c follows three steps: interpolation (similar to the regula falsi), truncation (adjusting the regula falsi similar to Regula falsi § Improvements in regula falsi) and then projection onto the minmax interval. The combination of these steps produces a simultaneously minmax optimal method with guarantees similar to interpolation based methods for smooth functions, and in practice will outperform both the bisection method and interpolation based methods applied to both smooth and non-smooth functions. Interpolation Many root-finding processes work by interpolation. This consists in using the last computed approximate values of the root for approximating the function by a polynomial of low degree, which takes the same values at these approximate roots. Then the root of the polynomial is computed and used as a new approximate value of the root of the function, and the process is iterated. Interpolating two values yields a line: a polynomial of degree one. This is the basis of the secant method. Regula falsi is also an interpolation method that interpolates two points at a time but it differs from the secant method by using two points that are not necessarily the last two computed points. Three values define a parabolic curve: a quadratic function. This is the basis of Muller's method. Iterative methods Although all root-finding algorithms proceed by iteration, an iterative root-finding method generally uses a specific type of iteration, consisting of defining an auxiliary function, which is applied to the last computed approximations of a root for getting a new approximation. The iteration stops when a fixed point of the auxiliary function is reached to the desired precision, i.e., when a new computed value is sufficiently close to the preceding ones. Newton's method (and similar derivative-based methods) Newton's method assumes the function f to have a continuous derivative. Newton's method may not converge if started too far away from a root. However, when it does converge, it is faster than the bisection method; its order of convergence is usually quadratic whereas the bisection method's is linear. Newton's method is also important because it readily generalizes to higher-dimensional problems. Householder's methods are a class of Newton-like methods with higher orders of convergence. The first one after Newton's method is Halley's method with cubic order of convergence. Secant method Replacing the derivative in Newton's method with a finite difference, we get the secant method. This method does not require the computation (nor the existence) of a derivative, but the price is slower convergence (the order of convergence is the golden ratio, approximately 1.62). A generalization of the secant method in higher dimensions is Broyden's method. Steffensen's method If we use a polynomial fit to remove the quadratic part of the finite difference used in the secant method, so that it better approximates the derivative, we obtain Steffensen's method, which has quadratic convergence, and whose behavior (both good and bad) is essentially the same as Newton's method but does not require a derivative. Fixed point iteration method We can use the fixed-point iteration to find the root of a function. Given a function which we have set to zero to find the root (), we rewrite the equation in terms of so that becomes (note, there are often many functions for each function). Next, we relabel each side of the equation as so that we can perform the iteration. Next, we pick a value for and perform the iteration until it converges towards a root of the function. If the iteration converges, it will converge to a root. The iteration will only converge if . As an example of converting to , if given the function , we will rewrite it as one of the following equations. , , , , or . Inverse interpolation The appearance of complex values in interpolation methods can be avoided by interpolating the inverse of f, resulting in the inverse quadratic interpolation method. Again, convergence is asymptotically faster than the secant method, but inverse quadratic interpolation often behaves poorly when the iterates are not close to the root. Combinations of methods Brent's method Brent's method is a combination of the bisection method, the secant method and inverse quadratic interpolation. At every iteration, Brent's method decides which method out of these three is likely to do best, and proceeds by doing a step according to that method. This gives a robust and fast method, which therefore enjoys considerable popularity. Ridders' method Ridders' method is a hybrid method that uses the value of function at the midpoint of the interval to perform an exponential interpolation to the root. This gives a fast convergence with a guaranteed convergence of at most twice the number of iterations as the bisection method. Roots of polynomials Finding roots in higher dimensions The bisection method has been generalized to higher dimensions; these methods are called generalized bisection methods. At each iteration, the domain is partitioned into two parts, and the algorithm decides - based on a small number of function evaluations - which of these two parts must contain a root. In one dimension, the criterion for decision is that the function has opposite signs. The main challenge in extending the method to multiple dimensions is to find a criterion that can be computed easily and guarantees the existence of a root. The Poincaré–Miranda theorem gives a criterion for the existence of a root in a rectangle, but it is hard to verify because it requires evaluating the function on the entire boundary of the rectangle. Another criterion is given by a theorem of Kronecker. It says that, if the topological degree of a function f on a rectangle is non-zero, then the rectangle must contain at least one root of f. This criterion is the basis for several root-finding methods, such as those of Stenger and Kearfott. However, computing the topological degree can be time-consuming. A third criterion is based on a characteristic polyhedron. This criterion is used by a method called Characteristic Bisection. It does not require computing the topological degree; it only requires computing the signs of function values. The number of required evaluations is at least , where D is the length of the longest edge of the characteristic polyhedron. Note that Vrahatis and Iordanidis prove a lower bound on the number of evaluations, and not an upper bound. A fourth method uses an intermediate value theorem on simplices. Again, no upper bound on the number of queries is given.	Mathematics	Real analysis	null
153353	https://en.wikipedia.org/wiki/Andromeda%20%28constellation%29	Andromeda (constellation)	Andromeda is one of the 48 constellations listed by the 2nd-century Greco-Roman astronomer Ptolemy, and one of the 88 modern constellations. Located in the northern celestial hemisphere, it is named for Andromeda, daughter of Cassiopeia, in the Greek myth, who was chained to a rock to be eaten by the sea monster Cetus. Andromeda is most prominent during autumn evenings in the Northern Hemisphere, along with several other constellations named for characters in the Perseus myth. Because of its northern declination, Andromeda is visible only north of 40° south latitude; for observers farther south, it lies below the horizon. It is one of the largest constellations, with an area of 722 square degrees. This is over 1,400 times the size of the full moon, 55% of the size of the largest constellation, Hydra, and over 10 times the size of the smallest constellation, Crux. Its brightest star, Alpheratz (Alpha Andromedae), is a binary star that has also been counted as a part of Pegasus, while Gamma Andromedae (Almach) is a colorful binary and a popular target for amateur astronomers. With a variable brightness similar to Alpheratz, Mirach (Beta Andromedae) is a red giant, its color visible to the naked eye. The constellation's most obvious deep-sky object is the naked-eye Andromeda Galaxy (M31, also called the Great Galaxy of Andromeda), the closest spiral galaxy to the Milky Way and one of the brightest Messier objects. Several fainter galaxies, including M31's companions M110 and M32, as well as the more distant NGC 891, lie within Andromeda. The Blue Snowball Nebula, a planetary nebula, is visible in a telescope as a blue circular object. In Chinese astronomy, the stars that make up Andromeda were members of four different constellations that had astrological and mythological significance; a constellation related to Andromeda also exists in Hindu mythology. Andromeda is the location of the radiant for the Andromedids, a weak meteor shower that occurs in November. History and mythology The uranography of Andromeda has its roots most firmly in the Greek tradition, though a female figure in Andromeda's location had appeared earlier in Babylonian astronomy. The stars that make up Pisces and the middle portion of modern Andromeda formed a constellation representing a fertility goddess, sometimes named as Anunitum or the Lady of the Heavens. Andromeda is known as "the Chained Lady" or "the Chained Woman" in English. It was known as Mulier Catenata ("chained woman") in Latin and al-Mar'at al Musalsalah in Arabic. It has also been called Persea ("Perseus's wife") or Cepheis ("Cepheus's daughter"), all names that refer to Andromeda's role in the Greco-Roman myth of Perseus, in which Cassiopeia, the queen of Aethiopia, bragged that her daughter was more beautiful than the Nereids, sea nymphs blessed with incredible beauty. Offended at her remark, the nymphs petitioned Poseidon to punish Cassiopeia for her insolence, which he did by commanding the sea monster Cetus to attack Aethiopia. Andromeda's panicked father, Cepheus, was told by the Oracle of Ammon that the only way to save his kingdom was to sacrifice his daughter to Cetus. She was chained to a rock by the sea but was saved by the hero Perseus, who in one version of the story used the head of Medusa to turn the monster into stone; in another version, by the Roman poet Ovid in his Metamorphoses, Perseus slew the monster with his diamond sword. Perseus and Andromeda then married; the myth recounts that the couple had nine children together – seven sons and two daughters – and founded Mycenae and its Persideae dynasty. After Andromeda's death Athena placed her in the sky as a constellation, to honor her. Three of the neighboring constellations (Perseus, Cassiopeia and Cepheus) represent characters in the Perseus myth, while Cetus retreats to beyond Pisces. It is connected with the constellation Pegasus. Andromeda was one of the original 48 constellations formulated by Ptolemy in his 2nd-century Almagest, in which it was defined as a specific pattern of stars. She is typically depicted with α Andromedae as her head, ο and λ Andromedae as her chains, and δ, π, μ, β, and γ her body and legs. However, there is no universal depiction of Andromeda and the stars used to represent her body, head, and chains. Arab astronomers were aware of Ptolemy's constellations, but they included a second constellation representing a fish overlapping Andromeda's body; the nose of this fish was marked by a hazy patch that ‍we ‍now ‍know ‍as ‍the ‍Andromeda Galaxy, ‍M31. Several stars from Andromeda and most of the stars in Lacerta were combined in 1787 by German astronomer Johann Bode to form Honores Friderici (also called Friedrichs Ehre). It was designed to honour King Frederick II of Prussia, but quickly fell into disuse. Since the time of Ptolemy, Andromeda has remained a constellation and is officially recognized by the International Astronomical Union. Like all those that date back to a pattern known to Ptolemy, it is attributed to a wider zone and thus many surrounding stars. In 1922, the IAU defined its recommended three-letter abbreviation, "And". The official boundaries of Andromeda were defined in 1930 by Belgian astronomer Eugène Delporte as a polygon of 36 segments. Its right ascension is between 22h 57.5m and 2h 39.3m and its declination is between 53.19° and 21.68° in the equatorial coordinate system. In non-Western astronomy In traditional Chinese astronomy, nine stars from Andromeda (including Beta Andromedae, Mu Andromedae, and Nu Andromedae), along with seven stars from Pisces, formed an elliptical constellation called "Legs" (奎宿). This constellation either represented the foot of a walking person or a wild boar. Gamma Andromedae and its neighbors were called "Teen Ta Tseang Keun" (天大将军, heaven's great general), representing honour in astrology and a great general in mythology. Alpha Andromedae and Gamma Pegasi together made "Wall" (壁宿), representing the eastern wall of the imperial palace and/or the emperor's personal library. For the Chinese, the northern swath of Andromeda formed a stable for changing horses (, 天厩, stable on sky) and the far western part, along with most of Lacerta, became Tengshe, a flying snake. An Arab constellation called "al-Hut" (the fish) was composed of several stars in Andromeda, M31, and several stars in Pisces. ν And, μ And, β And, η And, ζ And, ε And, δ And, π And, and 32 And were all included from Andromeda; ν Psc, φ Psc, χ Psc, and ψ1 Psc were included from Pisces. As per Hindu astronomy, Andromeda is known as Devyani Constellation while Cassiopeia is Sharmishta Constellation. Devyani and Sharmishta are wives of King Yayati (Perseus Constellation) who is the earliest patriarch of the Kuru and Yadu Clans that are mentioned frequently in epic Mahabharat. There is an interesting story of these three characters mentioned in Mahabharat. Devyani is the daughter of Guru Shukracharya while Shar. Hindu legends surrounding Andromeda are similar to the Greek myths. Ancient Sanskrit texts depict Antarmada chained to a rock, as in the Greek myth. Scholars believe that the Hindu and Greek astrological myths were closely linked; one piece of evidence cited is the similarity between the names "Antarmada" and "Andromeda". Andromeda is also associated with the Mesopotamian creation story of Tiamat, the goddess of Chaos. She bore many demons for her husband, Apsu, but eventually decided to destroy them in a war that ended when Marduk killed her. He used her body to create the constellations as markers of time for humans. In the Marshall Islands, Andromeda, Cassiopeia, Triangulum, and Aries are incorporated into a constellation representing a porpoise. Andromeda's bright stars are mostly in the body of the porpoise; Cassiopeia represents its tail and Aries its head. In the Tuamotu islands, Alpha Andromedae was called Takurua-e-te-tuki-hanga-ruki, meaning "Star of the wearisome toil", and Beta Andromedae was called Piringa-o-Tautu. Features Stars α And (Alpheratz, Sirrah) is the brightest star in this constellation. It is an A0p class binary star with an overall apparent visual magnitude of 2.1 and a luminosity of . It is 97 light-years from Earth. It represents Andromeda's head in Western mythology, however, the star's traditional Arabic names – Alpheratz and Sirrah, from the phrase surrat al-faras – sometimes translated as "navel of the steed". The Arabic names are a reference to the fact that α And forms an asterism known as the "Great Square of Pegasus" with 3 stars in Pegasus: α, β, and γ Peg. As such, the star was formerly considered to belong to both Andromeda and Pegasus, and was co-designated as "Delta Pegasi (δ Peg)", although this name is no longer formally used. β And (Mirach) is a red-hued giant star of type M0 located in an asterism known as the "girdle". It is 198 light-years away, has a magnitude of 2.06, and a luminosity of with a planet discovered orbiting this star (b). Its name comes from the Arabic phrase al-Maraqq meaning "the loins" or "the loincloth", a phrase translated from Ptolemy's writing. However, β And was mostly considered by the Arabs to be a part of al-Hut, a constellation representing a larger fish than Pisces at Andromeda's feet. γ And (Almach) is an orange-hued bright giant star of type K3 found at the southern tip of the constellation with an overall magnitude of 2.14. Almach is a multiple star with a yellow primary of magnitude 2.3 and a blue-green secondary of magnitude 5, separated by 9.7 arcseconds. British astronomer William Herschel said of the star: "[the] striking difference in the color of the 2 stars, suggests the idea of a sun and its planet, to which the contrast of their unequal size contributes not a little." The secondary, described by Herschel as a "fine light sky-blue, inclining to green", is itself a double star, with a secondary of magnitude 6.3 and a period of 61 years. The system is 358 light-years away. Almach was named for the Arabic phrase ʿAnaq al-Ard, which means "the earth-kid", an obtuse reference to an animal that aids a lion in finding prey. δ And is an orange-hued giant star of type K3 orange giant of magnitude 3.3. It is 105 light-years from Earth. ι And, κ, λ, ο, and ψ And form an asterism known as "Frederick's Glory", a name derived from a former constellation (Frederici Honores). ι And is a blue-white hued main-sequence star of type B8, 502 light-years from Earth; κ And is a white-hued main-sequence star of type B9 IVn, 168 light-years from Earth; λ And is a yellow-hued giant star of type G8, 86 light-years from Earth; ο And is a blue-white hued giant star of type B6, 679 light-years from Earth; and ψ And is a blue-white hued main-sequence star of type B7, 988 light-years from Earth. μ And is a white-hued main-sequence star of type A5 and magnitude 3.9. It is 130 light-years away. υ And (Titawin) is a magnitude 4.1 binary system that consists of one F-type dwarf and an M-type dwarf. The primary star has a planetary system with 4 confirmed planets, 0.96 times, 14.57 times, 10.19 times and 1.06 the mass of Jupiter. The system is 44 light-years from Earth. ξ And (Adhil) is a binary star 217 light-years away. The primary is an orange-hued giant star of type K0. π And is a blue-white hued binary star of magnitude 4.3 that is 598 light-years away. The primary is a main-sequence star of type B5. Its companion star is of magnitude 8.9. 51 And (Nembus) was assigned by Johann Bayer to Perseus, where he designated it "Upsilon Persei (υ Per)", but it was moved to Andromeda by the International Astronomical Union. It is 177 light-years from Earth and is an orange-hued giant star of type K3. 54 And was a former designation for φ Per. 56 And is an optical binary star. The primary is a yellow-hued giant star of type K0 with an apparent magnitude of 5.7 that is 316 light-years away. The secondary is an orange-hued giant star of type K0 and magnitude 5.9 that is 990 light-years from Earth. R And is a Mira-type variable star with a period of 409 days. Its maximum magnitude is 5.8 and its minimum magnitude is 14.8, and it is at a distance of 1,250 light-years. There are 6 other Mira variables in Andromeda. Z And is the M-type prototype for its class of variable stars. It ranges in magnitude from a minimum of 12.4 to a maximum of 8. It is 2,720 light-years away. Ross 248 (HH Andromedae) is the ninth-closest star to Earth at a distance of 10.3 light-years. It is a red-hued main-sequence BY Draconis variable star of type M6. 14 And (Veritate) is a yellow-hued giant star of type G8 that is 251 light-years away. It has a mass of and a radius of . It has one planet, 14 Andromedae b, discovered in 2008. It orbits at a distance of 0.83 astronomical units from its parent star every 186 days and has a mass of . Of the stars brighter than 4th magnitude (and those with measured luminosity), Andromeda has a relatively even distribution of evolved and main-sequence stars. Deep-sky objects Andromeda's borders contain many visible distant galaxies. The most famous deep-sky object in Andromeda is the spiral galaxy cataloged as Messier 31 (M31) or NGC 224 but known colloquially as the Andromeda Galaxy for the constellation. M31 is one of the most distant objects visible to the naked eye, 2.2 million light-years from Earth (estimates range up to 2.5 million light-years). It is seen under a dark, transparent sky as a hazy patch in the north of the constellation. M31 is the largest neighboring galaxy to the Milky Way and the largest member of the Local Group of galaxies. In absolute terms, M31 is approximately 200,000 light-years in diameter, twice the size of the Milky Way. It is an enormous – 192.4 by 62.2 arcminutes in apparent size – barred spiral galaxy similar in form to the Milky Way and at an approximate magnitude of 3.5, is one of the brightest deep-sky objects in the northern sky. Despite being visible to the naked eye, the "little cloud" near Andromeda's figure was not recorded until AD 964, when the Arab astronomer al-Sufi wrote his Book of Fixed Stars. M31 was first observed telescopically shortly after its invention, by Simon Marius in 1612. The future of the Andromeda and Milky Way galaxies may be interlinked: in about five billion years, the two could potentially begin an Andromeda–Milky Way collision that would spark extensive new star formation. American astronomer Edwin Hubble included M31 (then known as the Andromeda Nebula) in his groundbreaking 1923 research on galaxies. Using the 100-inch Hooker Telescope at Mount Wilson Observatory in California, he observed Cepheid variable stars in M31 during a search for novae, allowing him to determine their distance by using the stars as standard candles. The distance he found was far greater than the size of the Milky Way, which led him to the conclusion that many similar objects were "island universes" on their own. Hubble originally estimated that the Andromeda Galaxy was 900,000 light-years away, but Ernst Öpik's estimate in 1925 put the distance closer to 1.5 million light-years. The Andromeda Galaxy's two main companions, M32 and M110 (also known as NGC 221 and NGC 205, respectively) are faint elliptical galaxies that lie near it. M32, visible with a far smaller size of 8.7 by 6.4 arcminutes, compared to M110, appears superimposed on the larger galaxy in a telescopic view as a hazy smudge, M110 also appears slightly larger and distinct from the larger galaxy; M32 is 0.5° south of the core, M110 is 1° northwest of the core. M32 was discovered in 1749 by French astronomer Guillaume Le Gentil and has since been found to lie closer to Earth than the Andromeda Galaxy itself. It is viewable in binoculars from a dark site owing to its high surface brightness of 10.1 and overall magnitude of 9.0. M110 is classified as either a dwarf spheroidal galaxy or simply a generic elliptical galaxy. It is far fainter than M31 and M32, but larger than M32 with a surface brightness of 13.2, magnitude of 8.9, and size of 21.9 by 10.9 arcminutes. The Andromeda Galaxy has a total of 15 satellite galaxies, including M32 and M110. Nine of these lie in a plane, which has caused astronomers to infer that they have a common origin. These satellite galaxies, like the satellites of the Milky Way, tend to be older, gas-poor dwarf elliptical and dwarf spheroidal galaxies. Along with the Andromeda Galaxy and its companions, the constellation also features NGC 891 (Caldwell 23), a smaller galaxy just east of Almach. It is a barred spiral galaxy seen edge-on, with a dark dust lane visible down the middle. NGC 891 is incredibly faint and small despite its magnitude of 9.9, as its surface brightness of 14.6 indicates; it is 13.5 by 2.8 arcminutes in size. NGC 891 was discovered by the brother-and-sister team of William and Caroline Herschel in August 1783. This galaxy is at an approximate distance of 30 million light-years from Earth, calculated from its redshift of 0.002. Andromeda's most celebrated open cluster is NGC 752 (Caldwell 28) at an overall magnitude of 5.7. It is a loosely scattered cluster in the Milky Way that measures 49 arcminutes across and features approximately twelve bright stars, although more than 60 stars of approximately 9th magnitude become visible at low magnifications in a telescope. It is considered to be one of the more inconspicuous open clusters. The other open cluster in Andromeda is NGC 7686, which has a similar magnitude of 5.6 and is also a part of the Milky Way. It contains approximately 20 stars in a diameter of 15 arcminutes, making it a tighter cluster than NGC 752. There is one prominent planetary nebula in Andromeda: NGC 7662 (Caldwell 22). Lying approximately 3 degrees southwest of Iota Andromedae at a distance of about 4,000 light-years from Earth, the "Blue Snowball Nebula" is a popular target for amateur astronomers. It earned its popular name because it appears as a faint, round, blue-green object in a telescope, with an overall magnitude of 9.2. Upon further magnification, it is visible as a slightly elliptical annular disk that gets darker towards the center, with a magnitude 13.2 central star. The nebula has an overall magnitude of 9.2 and is 20 by 130 arcseconds in size. Meteor showers Each November, the Andromedids meteor shower appears to radiate from Andromeda. The shower peaks in mid-to-late November every year, but has a low peak rate of fewer than 2 meteors per hour. Astronomers have often associated the Andromedids with Biela's Comet, which was destroyed in the 19th century, but that connection is disputed. Andromedid meteors are known for being very slow and the shower itself is considered to be diffuse, as meteors can be seen coming from nearby constellations as well as from Andromeda itself. Andromedid meteors sometimes appear as red fireballs. The Andromedids were associated with the most spectacular meteor showers of the 19th century; the storms of 1872 and 1885 were estimated to have a peak rate of 2 meteors per second (a zenithal hourly rate of 10,000), prompting a Chinese astronomer to compare the meteors to falling rain. The Andromedids had another outburst on December 3–5, 2011, the most active shower since 1885, with a maximum zenithal hourly rate of 50 meteors per hour. The 2011 outburst was linked to ejecta from Comet Biela, which passed close to the Sun in 1649. None of the meteoroids observed were associated with material from the comet's 1846 disintegration. The observers of the 2011 outburst predicted outbursts in 2018, 2023, and 2036.	Physical sciences	Other	Astronomy
153520	https://en.wikipedia.org/wiki/Fish%20farming	Fish farming	Fish farming or pisciculture involves commercial breeding of fish, most often for food, in fish tanks or artificial enclosures such as fish ponds. It is a particular type of aquaculture, which is the controlled cultivation and harvesting of aquatic animals such as fish, crustaceans, molluscs and so on, in natural or pseudo-natural environments. A facility that releases juvenile fish into the wild for recreational fishing or to supplement a species' natural numbers is generally referred to as a fish hatchery. Worldwide, the most important fish species produced in fish farming are carp, catfish, salmon and tilapia. Global demand is increasing for dietary fish protein, which has resulted in widespread overfishing in wild fisheries, resulting in significant decrease in fish stocks and even complete depletion in some regions. Fish farming allows establishment of artificial fish colonies that are provided with sufficient feeding, protection from natural predators and competitive threats, access to veterinarian service, and easier harvesting when needed, while being separate from and thus do not usually impact the sustainable yields of wild fish populations. While fish farming is practised worldwide, China alone provides 62% of the world's farmed fish production. As of 2016, more than 50% of seafood was produced by aquaculture. In the last three decades, aquaculture has been the main driver of the increase in fisheries and aquaculture production, with an average growth of 5.3 percent per year in the period 2000–2018, reaching a record 82.1 million tonnes in 2018. Farming carnivorous fish such as salmon, however, does not always reduce pressure on wild fisheries, such farmed fish are usually fed fishmeal and fish oil extracted from wild forage fish. The 2008 global returns for fish farming recorded by the FAO totaled 33.8 million tonnes worth about US$60 billion. Although fish farming for food is the most widespread, another major fish farming industry provides living fish for the aquarium trade. The vast majority of freshwater fish in the aquarium trade originate from farms in Eastern and Southern Asia, eastern Europe, Florida and South America that use either indoor tank systems or outdoor pond systems, while farming of fish for the marine aquarium trade happens at a much smaller scale. In 2022 24% of fishers and fish farmers and 62% of workers in post-harvest sector were women. Major species Categories Aquaculture makes use of local photosynthetic production (extensive) or fish that are fed with external food supply (intensive). Extensive aquaculture Extensive aquaculture is the other form of fish farming. Extensive aquaculture is more basic than intensive aquaculture in that less effort is put into the husbandry of the fish. Extensive aquaculture is done in the ocean, natural and man-made lakes, bays, rivers, and Fiords. Fish are contained within these habitats by multiple mesh enclosures which also function as trapping nets during harvest (Figure 3) (4). Since fish are susceptible to the elements, site placement is essential to ensure the rapid growth of the targeted species. The drawback of these facilities is that they depend on the surrounding area for good water quality in order to reduce mortality and increase the survivorship and growth rate of the fish (19). Fish chosen for extensive aquaculture are very hardy and often do well in high densities. Seaweed, prawns, mussels, carp, tilapia, tuna and salmon are the most prominent forms of extensive aqua cultured seafood. Extensive aquaculture facilities have negative impacts on the environment as well. Natural habitats are destroyed in the development of man made ponds used for extensive aquaculture. In the Philippines, shrimp aquaculture is responsible for the destruction of thousands of acres of mangrove fields which serve as nurseries and living habitats for many marine organisms. Benthic habitats are being depleted due to the high amount of organic waste produced by the fish which settles below their pens(4). Phytoplankton and algae break down fecal matter and residual fish meal reducing the amount of available oxygen in the water column, which chokes and kills the Benthic organisms. Another serious problem acquainted with extensive aquaculture is the introduction of invasive species into ecosystems (10). Escaped fish increase the competition between organisms for limited resources. Also, when foreign fish interbreed with wild species, they upset the genetic variability of the species, making them more prone to disease and infection. The high density of fish in these mesh tanks is very tempting for predators of the sea and air (19). To protect the harvest from predators protective netting is set up at a high cost. Often times predatorial fish and mammals like seals, sharks, and tuna get caught in these barrier nets and die. Some farmers protect their stocks from predatorial birds such as pelicans and albatross by shooting these sometimes endangered creatures. Intensive aquaculture In these kinds of systems fish production per unit of surface can be increased at will, as long as sufficient oxygen, fresh water and food are provided. Because of the requirement of sufficient fresh water, a massive water purification system must be integrated in the fish farm. One way to achieve this is to combine hydroponic horticulture and water treatment, see below. The exception to this rule are cages which are placed in a river or sea, which supplements the fish crop with sufficient oxygenated water. Some environmentalists object to this practice. The cost of inputs per unit of fish weight is higher than in extensive farming, especially because of the high cost of fish feed. It must contain a much higher level of protein (up to 60%) than cattle feed and a balanced amino acid composition, as well. These higher protein-level requirements are a consequence of the higher feed efficiency of aquatic animals (higher feed conversion ratio [FCR], that is, kg of feed per kg of animal produced). Fish such as salmon have an FCR around 1.1 kg of feed per kg of salmon whereas chickens are in the 2.5 kg of feed per kg of chicken range. Fish do not use energy to keep warm, eliminating some carbohydrates and fats in the diet, required to provide this energy. This may be offset, though, by the lower land costs and the higher production which can be obtained due to the high level of input control. Aeration of the water is essential, as fish need a sufficient oxygen level for growth. This is achieved by bubbling, cascade flow, or aqueous oxygen. Catfish in genus Clarias can breathe atmospheric air and can tolerate much higher levels of pollutants than trout or salmon, which makes aeration and water purification less necessary and makes Clarias species especially suited for intensive fish production. In some Clarias farms, about 10% of the water volume can consist of fish biomass. The risk of infections by parasites such as fish lice, fungi (Saprolegnia spp.), intestinal worms (such as nematodes or trematodes), bacteria (e.g., Yersinia spp., Pseudomonas spp.), and protozoa (such as dinoflagellates) is similar to that in animal husbandry, especially at high population densities. However, animal husbandry is a larger and more technologically mature area of human agriculture and has developed better solutions to pathogen problems. Intensive aquaculture has to provide adequate water quality (oxygen, ammonia, nitrite, etc.) levels to minimize stress on the fish. This requirement makes control of the pathogen problem more difficult. Intensive aquaculture requires tight monitoring and a high level of expertise of the fish farmer. Very-high-intensity recycle aquaculture systems (RAS, also Recirculating Aquaculture Systems), where all the production parameters are controlled, are being used for high-value species. By recycling water, little is used per unit of production. However, the process has high capital and operating costs. The higher cost structures mean that RAS is economical only for high-value products, such as broodstock for egg production, fingerlings for net pen aquaculture operations, sturgeon production, research animals, and some special niche markets such as live fish. Raising ornamental coldwater fish (goldfish or koi), although theoretically much more profitable due to the higher income per weight of fish produced, has been successfully carried out only in the 21st century. The increased incidences of dangerous viral diseases of koi carp, together with the high value of the fish, has led to initiatives in closed-system koi breeding and growing in a number of countries. Today, a few commercially successful intensive koi-growing facilities are operating in the UK, Germany, and Israel. Some producers have adapted their intensive systems in an effort to provide consumers with fish that do not carry dormant forms of viruses and diseases. In 2016, juvenile Nile tilapia were given a food containing dried Schizochytrium in place of fish oil. When compared to a control group raised on regular food, they exhibited higher weight gain and better food-to-growth conversion, plus their flesh was higher in healthy omega-3 fatty acids. Fish farms Within intensive and extensive aquaculture methods, numerous specific types of fish farms are used; each has benefits and applications unique to its design. Cage system Fish cages are placed in lakes, bayous, ponds, rivers, or oceans to contain and protect fish until they can be harvested. The method is also called "off-shore cultivation" when the cages are placed in the sea. They can be constructed of a wide variety of components. Fish are stocked in cages, artificially fed, and harvested when they reach market size. A few advantages of fish farming with cages are that many types of waters can be used (rivers, lakes, filled quarries, etc.), many types of fish can be raised, and fish farming can co-exist with sport fishing and other water uses. Cage farming of fishes in open seas is also gaining in popularity. Given concerns of disease, poaching, poor water quality, etc., generally pond systems are considered simpler to start and easier to manage. Also, the past occurrences of cage-failures leading to escapes, have raised concern regarding the culture of non-native fish species in dam or open-water cages. On August 22, 2017, there was a massive failure of such cages at a commercial fishery in Washington state in Puget Sound, leading to release of nearly 300,000 Atlantic salmon in non-native waters. This is believed to risk endangering the native Pacific salmon species. Marine Scotland has kept records of caged fish escapes since 1999. They have recorded 357 fish escape incidents with 3,795,206 fish escaping into fresh and salt water. One company, Dawnfresh Farming Limited has been responsible for 40 incident and 152,790 Rainbow Trout escaping into freshwater lochs. Though the cage-industry has made numerous technological advances in cage construction in recent years, the risk of damage and escape due to storms is always a concern. Semi-submersible marine technology is beginning to impact fish farming. In 2018, 1.5 million salmon are in the middle of a year-long trial at Ocean Farm 1 off the coast of Norway. The semi-submersible project is the world's first deep-sea aquaculture project, and includes -high by -diameter pen made from a series of mesh-wire frames and nets. It is designed to disperse wastes better than more conventional farms in sheltered coastal waters, therefore supporting higher fish packing density. Copper-alloy nets Recently, copper alloys have become important netting materials in aquaculture. Copper alloys are antimicrobial, that is, they destroy bacteria, viruses, fungi, algae, and other microbes. In the marine environment, the antimicrobial/algaecidal properties of copper alloys prevent biofouling, which can briefly be described as the undesirable accumulation, adhesion, and growth of microorganisms, plants, algae, tube worms, barnacles, mollusks, and other organisms. The resistance of organism growth on copper alloy nets also provides a cleaner and healthier environment for farmed fish to grow and thrive. Traditional netting involves regular and labor-intensive cleaning. In addition to its antifouling benefits, copper netting has strong structural and corrosion-resistant properties in marine environments. Copper-zinc brass alloys are deployed in commercial-scale aquaculture operations in Asia, South America, and USA (Hawaii). Extensive research, including demonstrations and trials, are being implemented on two other copper alloys: copper-nickel and copper-silicon. Each of these alloy types has an inherent ability to reduce biofouling, cage waste, disease, and the need for antibiotics, while simultaneously maintaining water circulation and oxygen requirements. Other types of copper alloys are also being considered for research and development in aquaculture operations. In Southeast Asia, the traditional cage farming platform is called kelong. Open net pen system The open net pens system is a method that takes place in natural waters, such as rivers, lakes, near the coast or offshore. The breeders rear the fish in large cages floating in the water. The fish are living in natural water but are isolated with a net. Because the only barrier separating the fish from the surrounding environment is a net, this allows the water to flow from the ‘natural’ surrounding through the fish farms. The site of the fish farm is crucial for the farm to be a success or not. Before any fish farm is settled, it is highly recommended to be selective with the site location of the farm. The site must be examined on some essential elements. Important conditions on the location are: A good interchange of water and also a high replacement of bottom water. At all depths should be a good current condition. This is necessary because the organic particles should be able to be carried away using the current. A gravel and sand bottom are qualified for fish farming, although bottoms with silt and mud are not qualified. These should be avoided. A net should be at least or more above the bottom, so depth is important. Despite these important site conditions, the open net pen method was very popular in Norway and China. This is because of the cost friendliness and efficiency of this method. Negative external effects Because of the ocean's water flow and other reasons, open net pen culture is seen as a high-risk method for the environment. The flow allows chemicals, parasites, waste and diseases to spread in the enclosed environment, and this is not beneficial for the natural environment. Another negative consequence is the high escape rate of the cultured fish from these open net pens. These escaped fish also pose a high risk to the surrounding ecosystems. The amount of organic waste produced by fish farms is also alarming. A salmon farm in Scotland, for instance, is estimated to produce as much organic waste as equivalent to a town of people between 10,000 and 20,000 people each year. Today 50% of the world's seafood is farm-raised. Irrigation ditch or pond systems These use irrigation ditches or farm ponds to raise fish. The basic requirement is to have a ditch or pond that retains water, possibly with an above-ground irrigation system (many irrigation systems use buried pipes with headers). Using this method, water allotments can be stored in ponds or ditches, usually lined with bentonite clay. In small systems, the fish are often fed commercial fish food, and their waste products can help fertilize the fields. In larger ponds, the pond grows water plants and algae as fish food. Some of the most successful ponds grow introduced strains of plants, as well as introduced strains of fish. Control of water quality is crucial. Fertilizing, clarifying, and pH control of the water can increase yields substantially, as long as eutrophication is prevented and oxygen levels stay high. Yields can be low if the fish grow ill from electrolyte stress. Composite fish culture The composite fish culture system is a technology developed in India by the Indian Council of Agricultural Research in the 1970s. In this system, of both local and imported fish, a combination of five or six fish species is used in a single fish pond. These species are selected so that they do not compete for food among them by having different types of food habitats. As a result, the food available in all the parts of the pond is used. Fish used in this system include catla and silver carp (surface feeders), rohu (a column feeder), and mrigal and common carp (bottom feeders). Other fish also feed on the excreta of the common carp, and this helps contribute to the efficiency of the system which in optimal conditions produces 3000–6000 kg of fish per hectare per year. One problem with such composite fish culture is that many of these fish breed only during monsoon. Even if fish are collected from the wild, they can be mixed with other species, as well. Thus, a major problem in fish farming is the lack of availability of good-quality stock. To overcome this problem, ways have now been worked out to breed these fish in ponds using hormonal stimulation. This has ensured the supply of pure fish stock in desired quantities. Integrated recycling systems One of the largest problems with freshwater pisciculture is that it can use a million gallons of water per acre (about 1 m3 of water per m2) each year. Extended water purification systems allow for the reuse (recycling) of local water. The largest-scale pure fish farms use a system derived (admittedly much refined) from the New Alchemy Institute in the 1970s. Basically, large plastic fish tanks are placed in a greenhouse. A hydroponic bed is placed near, above or between them. When tilapia are raised in the tanks, they are able to eat algae, which naturally grow in the tanks when the tanks are properly fertilized. The tank water is slowly circulated to the hydroponic beds, where the tilapia waste feeds commercial plant crops. Carefully cultured microorganisms in the hydroponic bed convert ammonia to nitrates, and the plants are fertilized by the nitrates and phosphates.Other wastes are strained out by the hydroponic media, which double as an aerated pebble-bed filter. This system, properly tuned, produces more edible protein per unit area than any other. A wide variety of plants can grow well in the hydroponic beds. Most growers concentrate on herbs (e.g. parsley and basil), which command premium prices in small quantities all year long. The most common customers are restaurant wholesalers. Since the system lives in a greenhouse, it adapts to almost all temperate climates, and may also adapt to tropical climates. The main environmental impact is discharge of water that must be salted to maintain the fishes' electrolyte balance. Current growers use a variety of proprietary tricks to keep fish healthy, reducing their expenses for salt and wastewater discharge permits. Some veterinary authorities speculate that ultraviolet ozone disinfectant systems (widely used for ornamental fish) may play a prominent part in keeping the tilapia healthy with recirculated water. A number of large, well-capitalized ventures in this area have failed. Managing both the biology and markets is complicated. One future development is the combination of integrated recycling systems with urban farming as tried in Sweden by the Greenfish Initiative. Classic fry farming This is also called a "flow through system" . Trout and other sport fish are often raised from eggs to fry or fingerlings and then trucked to streams and released. Normally, the fry are raised in long, shallow, concrete tanks, fed with fresh stream water. The fry receive commercial fish food in pellets. While not as efficient as the New Alchemists' method, it is also far simpler and has been used for many years to stock streams with sport fish. European eel (Anguilla anguilla) aquaculturalists procure a limited supply of glass eels, juvenile stages of the European eel which swim north from the Sargasso Sea breeding grounds, for their farms. The European eel is threatened with extinction because of the excessive catch of glass eels by Spanish fishermen and overfishing of adult eels in, e.g., the Dutch IJsselmeer. Although European eel larvae can survive for several weeks, the full life cycle has not yet been achieved in captivity. Issues Welfare There is a growing consensus that fish can feel pain. Despite the vast number of fish consumed, fish welfare has historically received little attention. Farmed fish are usually raised in overcrowded environments, making them susceptible to stress, injuries, aggression and diseases. These conditions prevent them from engaging in natural behaviors such as nesting or migration. Overcrowding often leads to poor water quality due to fish waste and antibiotics use. Sea lice infestations are common and can cause painful lesion, but are typically treated with harsh chemicals. Additionally, fish are genetically engineered to grow larger and faster, leading to health problems such as cataracts and abnormal heart shapes. Feeding The issue of feeds in fish farming has been a controversial one. Many cultured fishes (tilapia, carp, catfish, many others) can be raised on a strictly herbivorous diet. Top-level carnivores (most salmonidae species in particular) on the other hand, depend on fish feed, of which a large portion is usually derived from wild-caught fish (anchovies, menhaden, etc.). Vegetable-derived proteins have successfully replaced fish meal in feeds for carnivorous fishes, but vegetable-derived oils have not successfully been incorporated into the diets of carnivores. Research is underway to try to change this, such that even salmon and other carnivores could be successfully fed with vegetable products. The F3 Challenge (Fish-Free Feed Challenge), as explained by a report from Wired in February 2017, "is a race to sell 100,000 metric tons of fish food, without the fish. Earlier this month, start-ups from places like Pakistan, China, and Belgium joined their American competition at the Google headquarters in Mountain View, California, showing off feed made from seaweed extracts, yeast, and algae grown in bioreactors." Not only do the feeds for carnivorous fish, like certain salmon species, remain controversial due to the containment of wild caught fish like anchovies, but they are not helping the health of the fish, as is the case in Norway. Between 2003 and 2007, Aldrin et al. examined three infectious diseases in Norwegian salmon fish farms—heart and skeletal muscle inflammation, pancreas disease, and infectious salmon anemia. In 2014, Martinez-Rubio et al. conducted a study in which cardiomyopathy syndrome (CMS), a severe cardiac disease in Atlantic salmon (Salmo salar), was investigated pertaining the effects of functional feeds with reduced lipid content and increased eicosapentaenoic acid levels in controlling CMS in salmon after infection with Piscine Myocarditis Virus (PMCV). Functional feeds are defined as high-quality feeds that beyond purposes of nutrition, they are formulated with health promoting features that could be beneficial in supporting disease resistance, such as CMS. Choosing a clinical nutrition approach using functional feeds could potentially move away from chemotherapeutic and antibiotic treatments, which could lower the costs of disease treatment and management in fish farms. In this investigation three fishmeal-based diets were served—one made of 31% lipid and the other two made of 18% lipid (one contained fishmeal and the other krill meal. Results demonstrated a significant difference in the immune and inflammatory responses and pathology in heart tissue as the fish were infected with PMCV. Fish fed with functional feeds with low lipid content demonstrated milder and delayed inflammatory response and therefore, less severe heart lesions at earlier and later stages after PMCV infection. Stocking density Secondly, farmed fish are kept in concentrations never seen in the wild (e.g. 50,000 fish in a area.). However, fish tend also to be animals that aggregate into large schools at high density. Most successful aquaculture species are schooling species, which do not have social problems at high density. Aquaculturists feel that operating a rearing system above its design capacity or above the social density limit of the fish will result in decreased growth rate and increased feed conversion ratio (kg dry feed/kg of fish produced), which results in increased cost and risk of health problems along with a decrease in profits. Stressing the animals is not desirable, but the concept of and measurement of stress must be viewed from the perspective of the animal using the scientific method. Parasites and disease Sea lice, particularly Lepeophtheirus salmonis and various Caligus species, including C. clemensi and C. rogercresseyi, can cause deadly infestations of both farm-grown and wild salmon. Sea lice are ectoparasites which feed on mucus, blood, and skin, and migrate and latch onto the skin of wild salmon during free-swimming, planktonic nauplii and copepodid larval stages, which can persist for several days. Large numbers of highly populated, open-net salmon farms can create exceptionally large concentrations of sea lice; when exposed in river estuaries containing large numbers of open-net farms, many young wild salmon are infected, and do not survive as a result. Adult salmon may survive otherwise critical numbers of sea lice, but small, thin-skinned juvenile salmon migrating to sea are highly vulnerable. On the Pacific coast of Canada, the louse-induced mortality of pink salmon in some regions is commonly over 80%. In Scotland, official figures show that more than nine million fish were lost to disease, parasites, botched treatment attempts and other problems on fish farms between 2016 and 2019. One of the treatments for parasite infestations involved bathing fish in hydrogen peroxide, which can harm or kill farmed fish if they are in a weak condition or if the chemical concentration is too strong. A 2008 meta-analysis of available data shows that salmon farming reduces the survival of associated wild salmon populations. This relationship has been shown to hold for Atlantic, steelhead, pink, chum, and coho salmon. The decrease in survival or abundance often exceeds 50%. Diseases and parasites are the most commonly cited reasons for such decreases. Some species of sea lice have been noted to target farmed coho and Atlantic salmon. Such parasites have been shown to have an effect on nearby wild fish. One place that has garnered international media attention is British Columbia's Broughton Archipelago. There, juvenile wild salmon must "run a gauntlet" of large fish farms located off-shore near river outlets before making their way to sea. The farms allegedly cause such severe sea lice infestations that one study predicted in 2007 a 99% collapse in the wild salmon population by 2011. This claim, however, has been criticized by numerous scientists who question the correlation between increased fish farming and increases in sea lice infestation among wild salmon. Because of parasite problems, some aquaculture operators frequently use strong antibiotic drugs to keep the fish alive, but many fish still die prematurely at rates up to 30%. Additionally, other common drugs used in salmonid fish farms in North America and Europe include anesthetic, chemotherapeutic, and anthelmintic agents. In some cases, these drugs have entered the environment. Additionally, the residual presence of these drugs in human food products has become controversial. Use of antibiotics in food production is thought to increase the prevalence of antibiotic resistance in human diseases. At some facilities, the use of antibiotic drugs in aquaculture has decreased considerably due to vaccinations and other techniques. However, most fish-farming operations still use antibiotics, many of which escape into the surrounding environment. The lice and pathogen problems of the 1990s facilitated the development of current treatment methods for sea lice and pathogens, which reduced the stress from parasite/pathogen problems. However, being in an ocean environment, the transfer of disease organisms from the wild fish to the aquaculture fish is an ever-present risk. Ecosystem impacts The large number of fish kept long-term in a single location contributes to habitat destruction of the nearby areas. The high concentrations of fish produce a significant amount of condensed faeces, often contaminated with drugs, which again affects local waterways. Aquaculture not only impacts the fish on the farm, but it also influences other species, which in return are attracted to or repelled by the farms. Mobile fauna, such as crustaceans, fish, birds, and marine mammals, interact with the process of aquaculture, but the long-term or ecological effects as a result of these interactions is still unknown. Some of these fauna may be attracted or demonstrate repulsion. The attraction/repulsion mechanism has various direct and indirect effects on wild organisms at individual and population levels. The interactions that wild organisms have with aquaculture may have implications on the management of fisheries species and the ecosystem in relation to how the fish farms are structured and organized. Siting If aquaculture farms are placed in an area with strong current, pollutants can be flushed out of the area fairly quickly. This helps manage the pollution problem and also aids in overall fish growth. Concern remains that resultant bacterial growth fertilised by fish faeces strips the water of oxygen, reducing or killing off the local marine life. Once an area has been so contaminated, fish farms are typically moved to new, uncontaminated areas. This practice has angered nearby fishermen. Other potential problems faced by aquaculturists include the obtaining of various permits and water-use rights, profitability, concerns about invasive species and genetic engineering depending on what species are involved, and interaction with the United Nations Convention on the Law of the Sea. Genetic engineering In regards to genetically engineered, farmed salmon, concern has been raised over their proven reproductive advantage and how it could potentially decimate local fish populations, if released into the wild. Biologist Rick Howard did a controlled laboratory study where wild fish and genetically engineered fish were allowed to breed. In 1989, AquaBounty Technologies developed the AquAdvantage salmon. The concerns and critiques of cultivating this genetically engineered fish in aquaculture are that the fish will escape and interact with other fish ultimately leading to the reproduction with other fishes. However, the FDA, has determined that while net pens would not be the most appropriate to prevent escapes, raising the salmon in Panama waters would effectively prevent escape because the water conditions there would fail to support long-term survival of any escaped salmon. Another method of preventing Aqua Advantage fish from impacting the ecosystems in the case they escape suggested by the FDA was to create sterile triploid females. This way concerns on reproducing with other fishes would be out of the question. The genetically engineered fish crowded out the wild fish in spawning beds, but the offspring were less likely to survive. The colorant used to make pen-raised salmon appear rosy like the wild fish has been linked with retinal problems in humans. Labeling In 2005, Alaska passed legislation requiring that any genetically altered fish sold in the state be labeled. In 2006, a Consumer Reports investigation revealed that farm-raised salmon is frequently sold as wild. In 2008, the US National Organic Standards Board allowed farmed fish to be labeled as organic provided less than 25% of their feed came from wild fish. This decision was criticized by the advocacy group Food & Water Watch as "bending the rules" about organic labeling. In the European Union, fish labeling as to species, method of production and origin has been required since 2002. Concerns continue over the labeling of salmon as farmed or wild-caught, as well as about the humane treatment of farmed fish. The Marine Stewardship Council has established an Eco label to distinguish between farmed and wild-caught salmon, while the RSPCA has established the Freedom Food label to indicate humane treatment of farmed salmon, as well as other food products. Indoor fish farming Other treatments such as ultraviolet sterilization, ozonation, and oxygen injection are also used to maintain optimal water quality. Through this system, many of the environmental drawbacks of aquaculture are minimized including escaped fish, water usage, and the introduction of pollutants. The practices also increased feed-use efficiency growth by providing optimum water quality. One of the drawbacks to recirculating aquaculture systems is the need for periodic water exchanges. However, the rate of water exchange can be reduced through aquaponics, such as the incorporation of hydroponically grown plants and denitrification. Both methods reduce the amount of nitrate in the water, and can potentially eliminate the need for water exchanges, closing the aquaculture system from the environment. The amount of interaction between the aquaculture system and the environment can be measured through the cumulative feed burden (CFB kg/M3), which measures the amount of feed that goes into the RAS relative to the amount of water and waste discharged. The environmental impact of larger indoor fish farming system will be linked to the local infrastructure, and water supply. Areas which are more drought-prone, indoor fish farms might flow out wastewater for watering agricultural farms, reducing the water affliction. From 2011, a team from the University of Waterloo led by Tahbit Chowdhury and Gordon Graff examined vertical RAS aquaculture designs aimed at producing protein-rich fish species. However, because of its high capital and operating costs, RAS has generally been restricted to practices such as broodstock maturation, larval rearing, fingerling production, research animal production, specific pathogen-free animal production, and caviar and ornamental fish production. As such, research and design work by Chowdhury and Graff remains difficult to implement. Although the use of RAS for other species is considered by many aquaculturalists to be currently impractical, some limited successful implementation of RAS has occurred with high-value product such as barramundi, sturgeon, and live tilapia in the US, eels and catfish in the Netherlands, trout in Denmark and salmon is planned in Scotland and Canada. Slaughter methods Tanks saturated with carbon dioxide have been used to make fish unconscious. Their gills are then cut with a knife so that the fish bleed out before they are further processed. This is no longer considered a humane method of slaughter. Methods that induce much less physiological stress are electrical or percussive stunning and this has led to the phasing out of the carbon dioxide slaughter method in Europe. Inhumane methods According to T. Håstein of the National Veterinary Institute (Oslo, Norway), "Different methods for slaughter of fish are in place and it is no doubt that many of them may be considered as appalling from an animal welfare point of view." A 2004 report by the EFSA Scientific Panel on Animal Health and Welfare explained: "Many existing commercial killing methods expose fish to substantial suffering over a prolonged period of time. For some species, existing methods, whilst capable of killing fish humanely, are not doing so because operators don't have the knowledge to evaluate them." Following are some less humane ways of killing fish. Air asphyxiation amounts to suffocation in the open air. The process can take upwards of 15 minutes to induce death, although unconsciousness typically sets in sooner. Ice baths or chilling of farmed fish on ice or submerged in near-freezing water is used to dampen muscle movements by the fish and to delay the onset of post-death decay. However, it does not necessarily reduce sensibility to pain; indeed, the chilling process has been shown to elevate cortisol. In addition, reduced body temperature extends the time before fish lose consciousness. CO narcosis Exsanguination without stunning is a process in which fish are taken up from water, held still, and cut so as to cause bleeding. According to references in Yue, this can leave fish writhing for an average of four minutes, and some catfish still responded to noxious stimuli after more than 15 minutes. Immersion in salt followed by gutting or other processing such as smoking is applied to eel. More humane methods Proper stunning renders the fish unconscious immediately and for a sufficient period of time such that the fish is killed in the slaughter process (e.g. through exsanguination) without regaining consciousness. Percussive stunning involves rendering the fish unconscious with a blow on the head. Electric stunning can be humane when a proper current is made to flow through the fish brain for a sufficient period of time. Electric stunning can be applied after the fish has been taken out of the water (dry stunning) or while the fish is still in the water. The latter generally requires a much higher current and may lead to operator safety issues. An advantage could be that in-water stunning allows fish to be rendered unconscious without stressful handling or displacement. However, improper stunning may not induce insensibility long enough to prevent the fish from enduring exsanguination while conscious. Whether the optimal stunning parameters that researchers have determined in studies are used by the industry in practice is unknown. Gallery	Technology	Aquaculture	null
153522	https://en.wikipedia.org/wiki/Plastid	Plastid	A plastid is a membrane-bound organelle found in the cells of plants, algae, and some other eukaryotic organisms. Plastids are considered to be intracellular endosymbiotic cyanobacteria. Examples of plastids include chloroplasts (used for photosynthesis); chromoplasts (used for synthesis and storage of pigments); leucoplasts (non-pigmented plastids, some of which can differentiate); and apicoplasts (non-photosynthetic plastids of apicomplexa derived from secondary endosymbiosis). A permanent primary endosymbiosis event occurred about 1.5 billion years ago in the Archaeplastida cladeland plants, red algae, green algae and glaucophytesprobably with a cyanobiont, a symbiotic cyanobacteria related to the genus Gloeomargarita. Another primary endosymbiosis event occurred later, between 140 to 90 million years ago, in the photosynthetic plastids Paulinella amoeboids of the cyanobacteria genera Prochlorococcus and Synechococcus, or the "PS-clade". Secondary and tertiary endosymbiosis events have also occurred in a wide variety of organisms; and some organisms developed the capacity to sequester ingested plastidsa process known as kleptoplasty. A. F. W. Schimper was the first to name, describe, and provide a clear definition of plastids, which possess a double-stranded DNA molecule that long has been thought of as circular in shape, like that of the circular chromosome of prokaryotic cellsbut now, perhaps not; (see "..a linear shape"). Plastids are sites for manufacturing and storing pigments and other important chemical compounds used by the cells of autotrophic eukaryotes. Some contain biological pigments such as used in photosynthesis or which determine a cell's color. Plastids in organisms that have lost their photosynthetic properties are highly useful for manufacturing molecules like the isoprenoids. In land plants Chloroplasts, proplastids, and differentiation In land plants, the plastids that contain chlorophyll can perform photosynthesis, thereby creating internal chemical energy from external sunlight energy while capturing carbon from Earth's atmosphere and furnishing the atmosphere with life-giving oxygen. These are the chlorophyll-plastidsand they are named chloroplasts; (see top graphic). Other plastids can synthesize fatty acids and terpenes, which may be used to produce energy or as raw material to synthesize other molecules. For example, plastid epidermal cells manufacture the components of the tissue system known as plant cuticle, including its epicuticular wax, from palmitic acidwhich itself is synthesized in the chloroplasts of the mesophyll tissue. Plastids function to store different components including starches, fats, and proteins. All plastids are derived from proplastids, which are present in the meristematic regions of the plant. Proplastids and young chloroplasts typically divide by binary fission, but more mature chloroplasts also have this capacity. Plant proplastids (undifferentiated plastids) may differentiate into several forms, depending upon which function they perform in the cell, (see top graphic). They may develop into any of the following variants: Chloroplasts: typically green plastids that perform photosynthesis. Etioplasts: precursors of chloroplasts. Chromoplasts: coloured plastids that synthesize and store pigments. Gerontoplasts: plastids that control the dismantling of the photosynthetic apparatus during plant senescence. Leucoplasts: colourless plastids that synthesize monoterpenes. Leucoplasts differentiate into even more specialized plastids, such as: the aleuroplasts; Amyloplasts: storing starch and detecting gravityfor maintaining geotropism. Elaioplasts: storing fats. Proteinoplasts: storing and modifying protein. or Tannosomes: synthesizing and producing tannins and polyphenols. Depending on their morphology and target function, plastids have the ability to differentiate or redifferentiate between these and other forms. Plastomes and Chloroplast DNA/ RNA; plastid DNA and plastid nucleoids Each plastid creates multiple copies of its own unique genome, or plastome, (from 'plastid genome')which for a chlorophyll plastid (or chloroplast) is equivalent to a 'chloroplast genome', or a 'chloroplast DNA'. The number of genome copies produced per plastid is variable, ranging from 1000 or more in rapidly dividing new cells, encompassing only a few plastids, down to 100 or less in mature cells, encompassing numerous plastids. A plastome typically contains a genome that encodes transfer ribonucleic acids (tRNA)s and ribosomal ribonucleic acids (rRNAs). It also contains proteins involved in photosynthesis and plastid gene transcription and translation. But these proteins represent only a small fraction of the total protein set-up necessary to build and maintain any particular type of plastid. Nuclear genes (in the cell nucleus of a plant) encode the vast majority of plastid proteins; and the expression of nuclear and plastid genes is co-regulated to coordinate the development and differention of plastids. Many plastids, particularly those responsible for photosynthesis, possess numerous internal membrane layers. Plastid DNA exists as protein-DNA complexes associated as localized regions within the plastid's inner envelope membrane; and these complexes are called 'plastid nucleoids'. Unlike the nucleus of a eukaryotic cell, a plastid nucleoid is not surrounded by a nuclear membrane. The region of each nucleoid may contain more than 10 copies of the plastid DNA. Where the proplastid (undifferentiated plastid) contains a single nucleoid region located near the centre of the proplastid, the developing (or differentiating) plastid has many nucleoids localized at the periphery of the plastid and bound to the inner envelope membrane. During the development/ differentiation of proplastids to chloroplastsand when plastids are differentiating from one type to anothernucleoids change in morphology, size, and location within the organelle. The remodelling of plastid nucleoids is believed to occur by modifications to the abundance of and the composition of nucleoid proteins. In normal plant cells long thin protuberances called stromules sometimes formextending from the plastid body into the cell cytosol while interconnecting several plastids. Proteins and smaller molecules can move around and through the stromules. Comparatively, in the laboratory, most cultured cellswhich are large compared to normal plant cellsproduce very long and abundant stromules that extend to the cell periphery. In 2014, evidence was found of the possible loss of plastid genome in Rafflesia lagascae, a non-photosynthetic parasitic flowering plant, and in Polytomella, a genus of non-photosynthetic green algae. Extensive searches for plastid genes in both taxons yielded no results, but concluding that their plastomes are entirely missing is still disputed. Some scientists argue that plastid genome loss is unlikely since even these non-photosynthetic plastids contain genes necessary to complete various biosynthetic pathways including heme biosynthesis. Even with any loss of plastid genome in Rafflesiaceae, the plastids still occur there as "shells" without DNA content, which is reminiscent of hydrogenosomes in various organisms. In algae and protists Plastid types in algae and protists include: Chloroplasts: found in green algae (plants) and other organisms that derived their genomes from green algae. Muroplasts: also known as cyanoplasts or cyanelles, the plastids of glaucophyte algae are similar to plant chloroplasts, excepting they have a peptidoglycan cell wall that is similar to that of bacteria. Rhodoplasts: the red plastids found in red algae, which allows them to photosynthesize down to marine depths of 268 m. The chloroplasts of plants differ from rhodoplasts in their ability to synthesize starch, which is stored in the form of granules within the plastids. In red algae, floridean starch is synthesized and stored outside the plastids in the cytosol. Secondary and tertiary plastids: from endosymbiosis of green algae and red algae. Leucoplast: in algae, the term is used for all unpigmented plastids. Their function differs from the leucoplasts of plants. Apicoplast: the non-photosynthetic plastids of Apicomplexa derived from secondary endosymbiosis. The plastid of photosynthetic Paulinella species is often referred to as the 'cyanelle' or chromatophore, and is used in photosynthesis. It had a much more recent endosymbiotic event, in the range of 140–90 million years ago, which is the only other known primary endosymbiosis event of cyanobacteria. Etioplasts, amyloplasts and chromoplasts are plant-specific and do not occur in algae. Plastids in algae and hornworts may also differ from plant plastids in that they contain pyrenoids . Inheritance In reproducing, most plants inherit their plastids from only one parent. In general, angiosperms inherit plastids from the female gamete, where many gymnosperms inherit plastids from the male pollen. Algae also inherit plastids from just one parent. Thus the plastid DNA of the other parent is completely lost. In normal intraspecific crossingsresulting in normal hybrids of one speciesthe inheriting of plastid DNA appears to be strictly uniparental; i.e., from the female. In interspecific hybridisations, however, the inheriting is apparently more erratic. Although plastids are inherited mainly from the female in interspecific hybridisations, there are many reports of hybrids of flowering plants producing plastids from the male. Approximately 20% of angiosperms, including alfalfa (Medicago sativa), normally show biparental inheriting of plastids. DNA damage and repair The plastid DNA of maize seedlings is subjected to increasing damage as the seedlings develop. The DNA damage is due to oxidative environments created by photo-oxidative reactions and photosynthetic/ respiratory electron transfer. Some DNA molecules are repaired but DNA with unrepaired damage is apparently degraded to non-functional fragments. DNA repair proteins are encoded by the cell's nuclear genome and then translocated to plastids where they maintain genome stability/ integrity by repairing the plastid's DNA. For example, in chloroplasts of the moss Physcomitrella patens, a protein employed in DNA mismatch repair (Msh1) interacts with proteins employed in recombinational repair (RecA and RecG) to maintain plastid genome stability. Origin Plastids are thought to be descended from endosymbiotic cyanobacteria. The primary endosymbiotic event of the Archaeplastida is hypothesized to have occurred around 1.5 billion years ago and enabled eukaryotes to carry out oxygenic photosynthesis. Three evolutionary lineages in the Archaeplastida have since emerged in which the plastids are named differently: chloroplasts in green algae and/or plants, rhodoplasts in red algae, and muroplasts in the glaucophytes. The plastids differ both in their pigmentation and in their ultrastructure. For example, chloroplasts in plants and green algae have lost all phycobilisomes, the light harvesting complexes found in cyanobacteria, red algae and glaucophytes, but instead contain stroma and grana thylakoids. The glaucocystophycean plastid—in contrast to chloroplasts and rhodoplasts—is still surrounded by the remains of the cyanobacterial cell wall. All these primary plastids are surrounded by two membranes. The plastid of photosynthetic Paulinella species is often referred to as the 'cyanelle' or chromatophore, and had a much more recent endosymbiotic event about 90–140 million years ago; it is the only known primary endosymbiosis event of cyanobacteria outside of the Archaeplastida. The plastid belongs to the "PS-clade" (of the cyanobacteria genera Prochlorococcus and Synechococcus), which is a different sister clade to the plastids belonging to the Archaeplastida. In contrast to primary plastids derived from primary endosymbiosis of a prokaryoctyic cyanobacteria, complex plastids originated by secondary endosymbiosis in which a eukaryotic organism engulfed another eukaryotic organism that contained a primary plastid. When a eukaryote engulfs a red or a green alga and retains the algal plastid, that plastid is typically surrounded by more than two membranes. In some cases these plastids may be reduced in their metabolic and/or photosynthetic capacity. Algae with complex plastids derived by secondary endosymbiosis of a red alga include the heterokonts, haptophytes, cryptomonads, and most dinoflagellates (= rhodoplasts). Those that endosymbiosed a green alga include the euglenids and chlorarachniophytes (= chloroplasts). The Apicomplexa, a phylum of obligate parasitic alveolates including the causative agents of malaria (Plasmodium spp.), toxoplasmosis (Toxoplasma gondii), and many other human or animal diseases also harbor a complex plastid (although this organelle has been lost in some apicomplexans, such as Cryptosporidium parvum, which causes cryptosporidiosis). The 'apicoplast' is no longer capable of photosynthesis, but is an essential organelle, and a promising target for antiparasitic drug development. Some dinoflagellates and sea slugs, in particular of the genus Elysia, take up algae as food and keep the plastid of the digested alga to profit from the photosynthesis; after a while, the plastids are also digested. This process is known as kleptoplasty, from the Greek, kleptes (), thief. Plastid development cycle In 1977 J.M Whatley proposed a plastid development cycle which said that plastid development is not always unidirectional but is instead a complicated cyclic process. Proplastids are the precursor of the more differentiated forms of plastids, as shown in the diagram to the right.	Biology and health sciences	Plant cells	null
153663	https://en.wikipedia.org/wiki/Cytokine	Cytokine	Cytokines () are a broad and loose category of small proteins (~5–25 kDa) important in cell signaling. Due to their size, cytokines cannot cross the lipid bilayer of cells to enter the cytoplasm and therefore typically exert their functions by interacting with specific cytokine receptors on the target cell surface. Cytokines have been shown to be involved in autocrine, paracrine and endocrine signaling as immunomodulating agents. Cytokines include chemokines, interferons, interleukins, lymphokines, and tumour necrosis factors, but generally not hormones or growth factors (despite some overlap in the terminology). Cytokines are produced by a broad range of cells, including immune cells like macrophages, B lymphocytes, T lymphocytes and mast cells, as well as endothelial cells, fibroblasts, and various stromal cells; a given cytokine may be produced by more than one type of cell. They act through cell surface receptors and are especially important in the immune system; cytokines modulate the balance between humoral and cell-based immune responses, and they regulate the maturation, growth, and responsiveness of particular cell populations. Some cytokines enhance or inhibit the action of other cytokines in complex ways. They are different from hormones, which are also important cell signaling molecules. Hormones circulate in higher concentrations, and tend to be made by specific kinds of cells. Cytokines are important in health and disease, specifically in host immune responses to infection, inflammation, trauma, sepsis, cancer, and reproduction. The word comes from the ancient Greek language: cyto, from Greek κύτος, kytos, 'cavity, cell' + kines, from Greek κίνησις, kinēsis, 'movement'. Discovery Interferon-alpha, an interferon type I, was identified in 1957 as a protein that interfered with viral replication. The activity of interferon-gamma (the sole member of the interferon type II class) was described in 1965; this was the first identified lymphocyte-derived mediator. Macrophage migration inhibitory factor (MIF) was identified simultaneously in 1966 by John David and Barry Bloom. In 1969, Dudley Dumonde proposed the term "lymphokine" to describe proteins secreted from lymphocytes and later, proteins derived from macrophages and monocytes in culture were called "monokines". In 1974, pathologist Stanley Cohen, M.D. (not to be confused with the Nobel laureate named Stanley Cohen, who was a PhD biochemist; nor with the MD geneticist Stanley Norman Cohen) published an article describing the production of MIF in virus-infected allantoic membrane and kidney cells, showing its production is not limited to immune cells. This led to his proposal of the term cytokine. In 1993, Ogawa described the early acting growth factors, intermediate acting growth factors and late acting growth factors. Difference from hormones Classic hormones circulate in aqueous solution in nanomolar (10 M) concentrations that usually vary by less than one order of magnitude. In contrast, some cytokines (such as IL-6) circulate in picomolar (10 M) concentrations that can increase up to 1,000 times during trauma or infection. The widespread distribution of cellular sources for cytokines may be a feature that differentiates them from hormones. Virtually all nucleated cells, but especially endo/epithelial cells and resident macrophages (many near the interface with the external environment) are potent producers of IL-1, IL-6, and TNF-α. In contrast, classic hormones, such as insulin, are secreted from discrete glands such as the pancreas. The current terminology refers to cytokines as immunomodulating agents. A contributing factor to the difficulty of distinguishing cytokines from hormones is that some immunomodulating effects of cytokines are systemic (i.e., affecting the whole organism) rather than local. For instance, to accurately utilize hormone terminology, cytokines may be autocrine or paracrine in nature, and chemotaxis, chemokinesis and endocrine as a pyrogen. Essentially, cytokines are not limited to their immunomodulatory status as molecules. Nomenclature Cytokines have been classed as lymphokines, interleukins, and chemokines, based on their presumed cell of secretion, function, or target of action. Because cytokines are characterised by considerable redundancy and pleiotropism, such distinctions, allowing for exceptions, are obsolete. The term interleukin was initially used by researchers for those cytokines whose presumed targets are principally white blood cells (leukocytes). It is now used largely for designation of newer cytokine molecules and bears little relation to their presumed function. The vast majority of these are produced by T-helper cells. Lymphokines: produced by lymphocytes Monokines: produced exclusively by monocytes Interferons: involved in antiviral responses Colony stimulating factors: support the growth of cells in semisolid media Chemokines: mediate chemoattraction (chemotaxis) between cells. Classification Structural Structural homogeneity has been able to partially distinguish between cytokines that do not demonstrate a considerable degree of redundancy so that they can be classified into four types: The four-α-helix bundle family (): member cytokines have three-dimensional structures with a bundle of four α-helices. This family, in turn, is divided into three sub-families: the IL-2 subfamily. This is the largest family. It contains several non-immunological cytokines including erythropoietin (EPO) and thrombopoietin (TPO). They can be grouped into long-chain and short-chain cytokines by topology. Some members share the common gamma chain as part of their receptor. the interferon (IFN) subfamily. the IL-10 subfamily. The IL-1 family, which primarily includes IL-1 and IL-18. The cysteine knot cytokines () include members of the transforming growth factor beta superfamily, including TGF-β1, TGF-β2 and TGF-β3. The IL-17 family, which has yet to be completely characterized, though member cytokines have a specific effect in promoting proliferation of T-cells that have cytotoxic effects. Functional A classification that proves more useful in clinical and experimental practice outside of structural biology divides immunological cytokines into those that enhance cellular immune responses, type 1 (TNFα, IFN-γ, etc.), and those that enhance antibody responses, type 2 (TGF-β, IL-4, IL-10, IL-13, etc.). A key focus of interest has been that cytokines in one of these two sub-sets tend to inhibit the effects of those in the other. Dysregulation of this tendency is under intensive study for its possible role in the pathogenesis of autoimmune disorders. Several inflammatory cytokines are induced by oxidative stress. The fact that cytokines themselves trigger the release of other cytokines and also lead to increased oxidative stress makes them important in chronic inflammation, as well as other immunoresponses, such as fever and acute phase proteins of the liver (IL-1,6,12, IFN-a). Cytokines also play a role in anti-inflammatory pathways and are a possible therapeutic treatment for pathological pain from inflammation or peripheral nerve injury. There are both pro-inflammatory and anti-inflammatory cytokines that regulate this pathway. Receptors In recent years, the cytokine receptors have come to demand the attention of more investigators than cytokines themselves, partly because of their remarkable characteristics and partly because a deficiency of cytokine receptors has now been directly linked to certain debilitating immunodeficiency states. In this regard, and also because the redundancy and pleomorphism of cytokines are, in fact, a consequence of their homologous receptors, many authorities think that a classification of cytokine receptors would be more clinically and experimentally useful. A classification of cytokine receptors based on their three-dimensional structure has, therefore, been attempted. Such a classification, though seemingly cumbersome, provides several unique perspectives for attractive pharmacotherapeutic targets. Immunoglobulin (Ig) superfamily, which are ubiquitously present throughout several cells and tissues of the vertebrate body, and share structural homology with immunoglobulins (antibodies), cell adhesion molecules, and even some cytokines. Examples: IL-1 receptor types. Hemopoietic Growth Factor (type 1) family, whose members have certain conserved motifs in their extracellular amino-acid domain. The IL-2 receptor belongs to this chain, whose γ-chain (common to several other cytokines) deficiency is directly responsible for the x-linked form of Severe Combined Immunodeficiency (X-SCID). Interferon (type 2) family, whose members are receptors for IFN β and γ. Tumor necrosis factors (TNF) (type 3) family, whose members share a cysteine-rich common extracellular binding domain, and includes several other non-cytokine ligands like CD40, CD27 and CD30, besides the ligands on which the family is named. Seven transmembrane helix family, the ubiquitous receptor type of the animal kingdom. All G protein-coupled receptors (for hormones and neurotransmitters) belong to this family. Chemokine receptors, two of which act as binding proteins for HIV (CD4 and CCR5), also belong to this family. Interleukin-17 receptor (IL-17R) family, which shows little homology with any other cytokine receptor family. Structural motifs conserved between members of this family include: an extracellular fibronectin III-like domain, a transmembrane domain and a cytoplasmic SERIF domain. The known members of this family are as follows: IL-17RA, IL-17RB, IL-17RC, IL17RD and IL-17RE. Cellular effects Each cytokine has a matching cell-surface receptor. Subsequent cascades of intracellular signaling then alter cell functions. This may include the upregulation and/or downregulation of several genes and their transcription factors, resulting in the production of other cytokines, an increase in the number of surface receptors for other molecules, or the suppression of their own effect by feedback inhibition. The effect of a particular cytokine on a given cell depends on the cytokine, its extracellular abundance, the presence and abundance of the complementary receptor on the cell surface, and downstream signals activated by receptor binding; these last two factors can vary by cell type. Cytokines are characterized by considerable redundancy, in that many cytokines appear to share similar functions. It seems to be a paradox that cytokines binding to antibodies have a stronger immune effect than the cytokine alone. This may lead to lower therapeutic doses. It has been shown that inflammatory cytokines cause an IL-10-dependent inhibition of T-cell expansion and function by up-regulating PD-1 levels on monocytes, which leads to IL-10 production by monocytes after binding of PD-1 by PD-L. Adverse reactions to cytokines are characterized by local inflammation and/or ulceration at the injection sites. Occasionally such reactions are seen with more widespread papular eruptions. Roles in health and disease Cytokines are involved in several developmental processes during embryonic development. Cytokines are released from the blastocyst, and are also expressed in the endometrium, and have critical roles in the stages of zona hatching, and implantation. Cytokines are crucial for fighting off infections and in other immune responses. However, they can become dysregulated and pathological in inflammation, trauma, sepsis, and hemorrhagic stroke. Dysregulated cytokine secretion in the aged population can lead to inflammaging, and render these individuals more vulnerable to age-related diseases like neurodegenerative diseases and type 2 diabetes. A 2019 review was inconclusive as to whether cytokines play any definitive role in ME/CFS. A 2024 study found a positive correlation between plasma interleukin IL-2 and fatigue in patients with type 1 narcolepsy. Adverse effects Adverse effects of cytokines have been linked to many disease states and conditions ranging from schizophrenia, major depression and Alzheimer's disease to cancer. T regulatory cells (Tregs) and related-cytokines are effectively engaged in the process of tumor immune escape and functionally inhibit immune response against the tumor. Forkhead box protein 3 (Foxp3) as a transcription factor is an essential molecular marker of Treg cells. Foxp3 polymorphism (rs3761548) might be involved in cancer progression like gastric cancer through influencing Tregs function and the secretion of immunomodulatory cytokines such as IL-10, IL-35, and TGF-β. Normal tissue integrity is preserved by feedback interactions between diverse cell types mediated by adhesion molecules and secreted cytokines; disruption of normal feedback mechanisms in cancer threatens tissue integrity. Over-secretion of cytokines can trigger a dangerous cytokine storm syndrome. Cytokine storms may have been the cause of severe adverse events during a clinical trial of TGN1412. Cytokine storms are also suspected to have been the main cause of death in the 1918 "Spanish Flu" pandemic. Deaths were weighted more heavily towards people with healthy immune systems, because of their ability to produce stronger immune responses, with dramatic increases in cytokine levels. Another example of cytokine storm is seen in acute pancreatitis. Cytokines are integral and implicated in all angles of the cascade, resulting in the systemic inflammatory response syndrome and multi-organ failure associated with this intra-abdominal catastrophe. In the COVID-19 pandemic, some deaths from COVID-19 have been attributable to cytokine release storms. Current data suggest cytokine storms may be the source of extensive lung tissue damage and dysfunctional coagulation in COVID-19 infections. Medical use as drugs Some cytokines have been developed into protein therapeutics using recombinant DNA technology. Recombinant cytokines being used as drugs as of 2014 include: Bone morphogenetic protein (BMP), used to treat bone-related conditions Erythropoietin (EPO), used to treat anemia Granulocyte colony-stimulating factor (G-CSF), used to treat neutropenia in cancer patients Granulocyte macrophage colony-stimulating factor (GM-CSF), used to treat neutropenia and fungal infections in cancer patients Interferon alfa, used to treat hepatitis C and multiple sclerosis Interferon beta, used to treat multiple sclerosis Interleukin 2 (IL-2), used to treat cancer. Interleukin 11 (IL-11), used to treat thrombocytopenia in cancer patients. Interferon gamma is used to treat chronic granulomatous disease and osteopetrosis	Biology and health sciences	Cell processes	Biology
153681	https://en.wikipedia.org/wiki/Celestial%20equator	Celestial equator	The celestial equator is the great circle of the imaginary celestial sphere on the same plane as the equator of Earth. By extension, it is also a plane of reference in the equatorial coordinate system. In other words, the celestial equator is an abstract projection of the terrestrial equator into outer space. Due to Earth's axial tilt, the celestial equator is currently inclined by about 23.44° with respect to the ecliptic (the plane of Earth's orbit), but has varied from about 22.0° to 24.5° over the past 5 million years due to perturbation from other planets. An observer standing on Earth's equator visualizes the celestial equator as a semicircle passing through the zenith, the point directly overhead. As the observer moves north (or south), the celestial equator tilts towards the opposite horizon. The celestial equator is defined to be infinitely distant (since it is on the celestial sphere); thus, the ends of the semicircle always intersect the horizon due east and due west, regardless of the observer's position on Earth. At the poles, the celestial equator coincides with the astronomical horizon. At all latitudes, the celestial equator is a uniform arc or circle because the observer is only finitely far from the plane of the celestial equator, but infinitely far from the celestial equator itself. Astronomical objects near the celestial equator appear above the horizon from most places on earth, but they culminate (reach the meridian) highest near the equator. The celestial equator currently passes through these constellations: These are the most globally visible constellations. Over thousands of years, the orientation of Earth's equator and thus the constellations the celestial equator passes through will change due to axial precession. Celestial bodies other than Earth also have similarly defined celestial equators.	Physical sciences	Celestial sphere: General	Astronomy
153761	https://en.wikipedia.org/wiki/Humber%20Bridge	Humber Bridge	The Humber Bridge is a single-span road suspension bridge near Kingston upon Hull, East Riding of Yorkshire, England. When it opened to traffic on 24 June 1981, it was the longest of its type in the world; the Akashi Kaikyō Bridge surpassed it in 1998, and it became the thirteenth-longest by 2024. The bridge spans the Humber (an estuary formed by the rivers Trent and Ouse), between Barton-upon-Humber on the south bank and Hessle on the north bank, connecting the East Riding of Yorkshire with North Lincolnshire. Both sides of the bridge were in the non-metropolitan county of Humberside until its dissolution in 1996. The bridge can be seen for miles around, from as far as Patrington in the East Riding of Yorkshire, and from out to sea miles off the coast. It is a Grade I listed building. By 2006, the bridge carried an average of 120,000 vehicles per week. The toll was £3.00 each way for cars (higher for commercial vehicles), which made it the most expensive toll crossing in the United Kingdom. In April 2012, the toll was halved to £1.50 each way after the UK government deferred £150 million from the bridge's outstanding debt. History Before the bridge, commuters crossed the Humber on the Humber Ferry from Corporation Pier at Hull and New Holland Pier at New Holland, Lincolnshire, or by road via the M62 (from 1976), M18 (from 1979) and M180 motorways, crossing, by way of the Ouse Bridge, the River Ouse near Goole (connected to the Humber). Until the mid-1970s the route south was via the single-carriageway A63 and the A614 (via grid-locked Thorne) where it met the busy A18 and crossed the Stainforth and Keadby Canal at Keadby Bridge, a swing bridge, which formed a bottleneck on the route, and on through Finningley and Bawtry, meeting the east–west A631. The journey was along straight single-carriageway roads across foggy moors interrupted by bottlenecks for most of the journey to Blyth, Nottinghamshire, where it met the A1, and the accident rate was high. Debates in Parliament were held on the low standard of the route across the windswept plains around Goole. It was not unexpected that under these conditions, a Humber Bridge, with connecting dual-carriageway approach roads and grade-separated junctions, would seem worthwhile. By the time the bridge opened, much of this inferior route had been transformed by dualling of the A63 and its bypasses, extending the M62 and the connecting of the M18 from Thorne to Wadworth. The obvious need for a Humber Bridge had been reduced by the late 1970s with the improvements of the motorway infrastructure in the region. Although welcome, these improvements detracted from the need for vehicles to cross a bridge from Hessle to Barton. The Humber Bridge was a victim of the success of the M62 before it opened. A hovercraft service, Minerva and Mercury, linked Hull Pier and Grimsby Docks from February to October 1969 but suffered relatively frequent breakdowns. Act of Parliament Plans for a bridge were drawn up in the 1930s when a team of engineers compiled a report on whether to bridge or tunnel the estuary. It was decided that a bridge would cost £1,750,000 over a tunnel which was costed at £7,200,000. Revised plans were unveiled in 1955, but work did not begin until 27 July 1972. The (7 & 8 Eliz. 2. c. xlvi), was promoted by Kingston Upon Hull Corporation and established the Humber Bridge Board to manage and raise funds to build the bridge and buy the land required for the approach roads. 1966 Kingston upon Hull North by-election The allocation of funds proved impossible until the 1966 Kingston upon Hull North by-election. Labour Prime Minister Harold Wilson prevailed upon his Minister of Transport Barbara Castle to sanction the building of the bridge. Dismay at the long wait for a crossing led to Christopher Rowe writing a protest song, "The Humber Bridge". Design The consulting engineers for the project were Freeman Fox & Partners (now Arcadis NV). Sir Ralph Freeman had produced the first ideas in 1927 and in the early 1930s the cost of the project was estimated at £1.725 million and that the bridge would be unlikely to recoup the construction or maintenance costs. In 1935 he had an idea for a suspension bridge for the Humber Tunnel Executive Committee. Sir Gilbert Roberts produced more ideas in 1955 for a bridge with a central span, costing £15 million, to be paid for by East Riding County Council and Lindsey County Council. When it became likely that a bridge would be constructed, Imperial College-educated Bernard Wex OBE (1922–1990) produced the design in 1964 that was actually built. The bridge was built to last 120 years. In 1985 Wex was awarded the Telford Medal by the Institution of Civil Engineers. In the 1950s he had helped to design High Marnham Power Station. He was a former UK chairman of the International Association of Bridge and Structural Engineers and helped to found the Steel Construction Institute in 1976. The architect was R. E. Slater ARIBA. The administration building (for the tolls) was designed by Parker & Rosner. The landscaping was designed by Prof Arnold Weddle. Wind tunnel testing took place at the National Maritime Institute (now part of BMT Limited) at Teddington, and the road deck was designed for wind speeds up to , but storms featuring considerably lower wind speeds have been cited as grounds for emergency repairs in recent years. Construction The main contractor for the steel superstructure was British Bridge Builders (the same grouping as for the Forth and Severn Road Bridges comprising Sir William Arrol & Co., then a unit of NEI Cranes Ltd, Cleveland Bridge & Engineering Company, and Redpath Dorman Long Ltd). The contractor for the concrete towers, anchorages and sub-structure was John Howard & Co Ltd of Chatham, Kent, which was later bought by Amec. Concrete was chosen for the towers, instead of steel, partly due to cost, but also to suit the landscape. Work began on the southern approach road in July 1972 by Clugston Construction of Scunthorpe. The approach road to the A1077 junction, by Costain Civil Engineering, began in September 1976. It included a span from the southern anchorage of seven pre-stressed concrete box sections and the A1077 junction, costing £4.25 million. Work on the bridge substructure (foundations) began in March 1973. To reduce heat of hydration in the concrete, which produces calcium silicate hydrate from belite, as much as 60 per cent of the Portland cement was replaced with ground granulated blast-furnace slag (GGBS). It took longer to build the southern anchorage due to a diaphragm wall design due to there not being enough shallow bedrock. The main southern approach roads from Barton to the M180 motorway junction at Barnetby were built in the late 1970s by Clugston Construction of Scunthorpe, opening in 1978. The towers were constructed by slipforming and the north tower was completed by May 1974. The southern foundations were completed in September 1975, with the pier completed in March 1976, and the southern tower was completed by September 1976; the bridge had been planned to open in 1976. The northern tower and anchorage was built on solid chalk but the southern tower and anchorage were built on fissured Kimmeridge Clay, from the southern shore and built with a difficult caisson design. The subcontractor for the concrete was Tileman & Co. of Shipston-on-Stour, south Warwickshire. Cable spinning took place between September 1977 and July 1979. Each cable weighs , with 37 strands of 404 lengths of cable. The cable on the northern span has four extra strands. Each cable can take a load of . The deck is of box girder form, the box sections around each. The first box sections were assembled in June 1975 and put into the main span on 9 November 1979. The toll buildings and north approach road were built by A. F. Budge of Retford, Nottinghamshire, costing £2.9 million. Work began on the administration building in November 1976. The toll system was manufactured by Plessey Controls of Poole, Dorset. Corrosion resistance on the steelwork was provided by Camrex Corrosion of Bellshill, North Lanarkshire. The road was laid by Tarmac Roadstone of Wolverhampton with mastic asphalt. In 2017, the bridge was designated a Grade I listed building. A-frames At road level the deck was fastened to the towers through four rocking A-frames, to allow for movement caused by the catenary supporting the deck from above deflecting with the weight of passing traffic, from thermal expansion, and from changes in wind loading. The devices catered for a maximum deflection of 2 metres. By 2011 it was noticed that the pivot-pin bearings carrying the frames had worn, allowing them to drop towards the support structure. Each frame was replaced by two new components: a vertical linkage to cater for longitudinal movement and a sliding bearing for lateral displacement. Opening The bridge opened to traffic on 24 June 1981 at a final cost of £91 million (). It was opened officially by Queen Elizabeth II on 17 July 1981, in a ceremony that included a prayer of dedication by the Archbishop of York and a fly-past by the Red Arrows. World record With a centre span of and a total length of , the Humber Bridge was the longest single-span suspension bridge in the world for 17 years, until the Akashi Kaikyō Bridge opened in Japan on 5 April 1998. Local benefits The road-distance between Hull and Grimsby fell by nearly ; the town of Scunthorpe and environs were relieved of the passing traffic between Hull and Grimsby. Bridge statistics The bridge's surface takes the form of a dual carriageway with a lower-level foot and cycle path on both sides. There is a permanent speed limit on the full length of the bridge. Each tower consists of a pair of hollow vertical concrete columns, each tall and tapering from square at the base to at the top. The bridge is designed to tolerate constant motion and bends more than in winds of . The towers, although vertical, are farther apart at the top than the bottom due to the curvature of the Earth. The total length of the suspension cable is . The north tower is on the bank and has foundations down to . The south tower is in the water, and descends to as a consequence of the shifting sandbanks that make up the estuary. The bridge held the record for the world's longest single-span suspension bridge for 17 years, from its opening in July 1981 until the opening of the Akashi Kaikyō Bridge in April 1998. In June 2024, it became the thirteenth-longest, single-span suspension bridge. The central span, at , is the longest in Britain and in the Western Hemisphere. It remains the longest single-span suspension bridge in the world that can be crossed on foot or by bicycle. The bridge is crossed twice during the annual Humber Bridge Half Marathon in June, and Hull Marathon in September. Toll update project In July 2013, work began on introducing a new electronic tolling system. The existing Humber Bridge toll system was largely the same as it was when the bridge opened in 1981. The computer system was over 15 years old, absorbed an increasing amount of maintenance, and needed to be replaced. The project would decrease waiting times and was welcomed by business and transport leaders. In the first phase, the toll booths and the toll plaza canopy were replaced, and in the second phase, writing, testing and setting up the new toll system was completed. From 2015 bridge users could set up an account with the bridge and pay into it. Account holders receive a device called the HumberTAG, a small electronic tag that enables the system to recognise the bridge user; the toll is automatically deducted from the user's account. Two central lanes through the plaza are free-flowing; they do not have booths and account holders are able to cross the bridge without stopping. Incidents and suicides During construction of the bridge, the road deck sections were floated up on barges then hoisted into place by cables. During one of these lifting operations some of the cables on two of the road deck sections failed, leaving the sections hanging at an angle. The sections were subsequently installed. On more than 200 occasions, people have jumped or fallen from the bridge since it was opened in 1981; only five people have survived. Between 1990 and February 2001 the Humber Rescue Team launched its boat 64 times to deal with people falling or jumping off the bridge. Notable incidents include the cases of a West Yorkshire woman and her two-year-old daughter who fell off the bridge in 2005 and that of a man jumping from the bridge to his death on the A63 road below in September 2006. Plans were announced on 26 December 2009 to construct a suicide barrier along the walkways of the bridge; design constraints were cited as the reason for not installing barriers during the construction of the bridge. In May 2017, a YouTuber with the username 'Night Scape', along with a small group, illegally scaled the bridge without safety equipment. The group of young men climbed up the structure to the top of the bridge using the suspension wires as handholds. Humberside Police and the Humber Bridge Board are reviewing the security measures. On 3 April 2021, the Humber Bridge was closed to pedestrians and cyclists following an unspecified 'recent incident'. The decision came after multiple deaths at the bridge in the month of March. Following the death of one individual that month, a petition calling for increased safety measures to 'secure' the bridge had gained thousands of signatures. Concerns were raised over how the change will affect those who commute on foot or by bike. On 6 May 2021, the bridge was reopened to pedestrians and cyclists between the hours of 0500 and 2100; only users registered in advance could use the bridge outside of those hours. More CCTV and notices were erected and more staff assigned to patrol the crossing. Finances The bridge had a toll charge of £1.50 for cars from 1 April 2012, until for six months it was £3.00 and the only trunk road British toll bridge to charge motorcycles (£1.20). In 2004 many motorcyclists held a slow-pay protest, taking off gloves and helmets and paying the toll in large denomination bank notes. Police reported that the protest caused a queue long. In 1996, Parliament passed the Humber Bridge (Debts) Act 1996 to reorganise the board's liabilities to ensure the bridge could be safely maintained. Much of the interest on the debt was suspended and deferred in a refinancing which saw no write off – the balance was to be paid using tolls. In 2006, Shona McIsaac, Labour MP for Cleethorpes, tabled a Private Member's Bill, the Humber Bridge Bill. The bill would have made amendments to the Humber Bridge Act 1959 (7 & 8 Eliz. 2. c. xlvi) "requiring the Secretary of State to give directions to members of the Humber Bridge Board regarding healthcare and to review the possibility of facilitating journeys across the Humber Bridge in relation to healthcare". The aim was to allow patients travelling between the banks for medical treatment to cross without paying the toll and to allow the Secretary of State for Transport to appoint two members of the board to represent the interests of the NHS. Even though the Bill received cross-party support (it was co-sponsored by Shadow Home Secretary David Davis and supported by all other MPs representing North Lincolnshire and the East Riding of Yorkshire) it ran out of time later that year. A protest at the bridge on 1 September 2007 was supported by the local Cancer Patients Involvement Group, the Road Haulage Association, Diana Wallis (MEP for Yorkshire and the Humber) and local business and council representatives. The government responded to the petition on 14 January 2008, stating that "Concessions or exemptions from tolls on the Humber Bridge are a matter for the Humber Bridge Board". In October 2008, a joint campaign was launched by the Scunthorpe Telegraph, Hull Daily Mail and Grimsby Telegraph to abolish the toll. The papers' campaign, A Toll Too Far, was launched after a mooted increase in the toll, receiving much support from councillors and MPs serving Lincolnshire and the East Riding of Yorkshire. The campaign was to stave off a potential increase, secure a reduction to £1.00 and ultimately to be abolished. Thousands of readers backed the campaign with a paper and an online petition. A public inquiry into the tolls was held in March 2009 by independent inspector Neil Taylor. In July 2009, the Department for Transport announced that it had decided not to allow the proposed increase. Transport Minister Sadiq Khan said he did not believe it was right for the tolls to be raised in the current economic climate. In October 2009, the government approved a £6 million grant for maintenance costs, which meant that there would be no toll increase before 2011 at the earliest, by which time tolls would have been frozen for five years. The board applied again to the Department of Transport in September 2010, to raise the tolls from April 2011 but the government ordered a public inquiry into the application. A three-day public inquiry was held in Hull in early March 2011. Following the recommendation by the planning inspector, the government gave approval, on 14 June 2011, for the increase. The toll was raised on 1 October 2011, at which point it became the most expensive in the United Kingdom. The Severn Bridge/Second Severn Crossing charged £5.70 for Wales-bound traffic. In the 2011 Autumn Statement on 29 November, the Chancellor of the Exchequer, George Osborne, announced that the government had agreed to reduce the debt on the bridge by £150 million, which would allow the toll for cars to be halved to £1.50. Following the government accepting the agreement, between the four local councils, to meet a portion of the debt if revenues proved insufficient, the Transport Secretary, Justine Greening, confirmed the reduction on 29 February 2012, with effect from April. Image gallery	Technology	Bridges	null
153783	https://en.wikipedia.org/wiki/Crystal%20optics	Crystal optics	Crystal optics is the branch of optics that describes the behaviour of light in anisotropic media, that is, media (such as crystals) in which light behaves differently depending on which direction the light is propagating. The index of refraction depends on both composition and crystal structure and can be calculated using the Gladstone–Dale relation. Crystals are often naturally anisotropic, and in some media (such as liquid crystals) it is possible to induce anisotropy by applying an external electric field. Isotropic media Typical transparent media such as glasses are isotropic, which means that light behaves the same way no matter which direction it is travelling in the medium. In terms of Maxwell's equations in a dielectric, this gives a relationship between the electric displacement field D and the electric field E: where ε0 is the permittivity of free space and P is the electric polarization (the vector field corresponding to electric dipole moments present in the medium). Physically, the polarization field can be regarded as the response of the medium to the electric field of the light. Electric susceptibility In an isotropic and linear medium, this polarization field P is proportional and parallel to the electric field E: where χ is the electric susceptibility of the medium. The relation between D and E is thus: where is the dielectric constant of the medium. The value 1+χ is called the relative permittivity of the medium, and is related to the refractive index n, for non-magnetic media, by Anisotropic media In an anisotropic medium, such as a crystal, the polarisation field P is not necessarily aligned with the electric field of the light E. In a physical picture, this can be thought of as the dipoles induced in the medium by the electric field having certain preferred directions, related to the physical structure of the crystal. This can be written as: Here χ is not a number as before but a tensor of rank 2, the electric susceptibility tensor. In terms of components in 3 dimensions: or using the summation convention: Since χ is a tensor, P is not necessarily colinear with E. In nonmagnetic and transparent materials, χij = χji, i.e. the χ tensor is real and symmetric. In accordance with the spectral theorem, it is thus possible to diagonalise the tensor by choosing the appropriate set of coordinate axes, zeroing all components of the tensor except χxx, χyy and χzz. This gives the set of relations: The directions x, y and z are in this case known as the principal axes of the medium. Note that these axes will be orthogonal if all entries in the χ tensor are real, corresponding to a case in which the refractive index is real in all directions. It follows that D and E are also related by a tensor: Here ε is known as the relative permittivity tensor or dielectric tensor. Consequently, the refractive index of the medium must also be a tensor. Consider a light wave propagating along the z principal axis polarised such the electric field of the wave is parallel to the x-axis. The wave experiences a susceptibility χxx and a permittivity εxx. The refractive index is thus: For a wave polarised in the y direction: Thus these waves will see two different refractive indices and travel at different speeds. This phenomenon is known as birefringence and occurs in some common crystals such as calcite and quartz. If χxx = χyy ≠ χzz, the crystal is known as uniaxial. (See Optic axis of a crystal.) If χxx ≠ χyy and χyy ≠ χzz the crystal is called biaxial. A uniaxial crystal exhibits two refractive indices, an "ordinary" index (no) for light polarised in the x or y directions, and an "extraordinary" index (ne) for polarisation in the z direction. A uniaxial crystal is "positive" if ne > no and "negative" if ne < no. Light polarised at some angle to the axes will experience a different phase velocity for different polarization components, and cannot be described by a single index of refraction. This is often depicted as an index ellipsoid. Other effects Certain nonlinear optical phenomena such as the electro-optic effect cause a variation of a medium's permittivity tensor when an external electric field is applied, proportional (to lowest order) to the strength of the field. This causes a rotation of the principal axes of the medium and alters the behaviour of light travelling through it; the effect can be used to produce light modulators. In response to a magnetic field, some materials can have a dielectric tensor that is complex-Hermitian; this is called a gyro-magnetic or magneto-optic effect. In this case, the principal axes are complex-valued vectors, corresponding to elliptically polarized light, and time-reversal symmetry can be broken. This can be used to design optical isolators, for example. A dielectric tensor that is not Hermitian gives rise to complex eigenvalues, which corresponds to a material with gain or absorption at a particular frequency.	Physical sciences	Crystallography	Physics
153784	https://en.wikipedia.org/wiki/Naginata	Naginata	The naginata (, ) is a polearm and one of several varieties of traditionally made Japanese blades (nihontō). Naginata were originally used by the samurai class of feudal Japan, as well as by ashigaru (foot soldiers) and sōhei (warrior monks). The naginata is the iconic weapon of the onna-musha, a type of female warrior belonging to the Japanese nobility. A common misconception is that the Naginata is a type of sword, rather than a polearm. Description A naginata consists of a wooden or metal pole with a curved single-edged blade on the end; it is similar to the Chinese guan dao or the European glaive. Similar to the katana, naginata often have a round handguard (tsuba) between the blade and shaft, when mounted in a koshirae (furniture). The 30 cm to 60 cm (11.8 inches to 23.6 inches) naginata blade is forged in the same manner as traditional Japanese swords. The blade has a long tang (nakago) which is inserted in the shaft. The blade is removable and is secured by means of a wooden peg called mekugi (目釘) that passes through a hole (mekugi-ana) in both the tang and the shaft. The shaft ranges from 120 cm to 240 cm (47.2 inches to 94.5 inches) in length and is oval shaped. The area of the shaft where the tang sits is the tachiuchi or tachiuke. The tachiuchi/tachiuke would be reinforced with metal rings (naginata dogane or semegane), and/or metal sleeves (sakawa) and wrapped with cord (san-dan maki). The end of the shaft has a heavy metal end cap (ishizuki or hirumaki). When not in use the blade would be covered with a wooden sheath. History It is assumed that the naginata was developed from an earlier weapon type of the later 1st millennium AD, the hoko yari. Another assumption is that the naginata was developed by lengthening the hilt of the tachi at the end of the Heian period, and it is not certain which theory is correct. It is generally believed that naginata first appeared in the Heian period (794–1185). The term naginata first appeared in historical documents in the Heian period. The earliest clear references to naginata date from 1146. In Honchō Seiki compiled from 1150 to 1159 in the late Heian period, it is recorded that Minamoto no Tsunemitsu mentioned that his weapon was a naginata. In the early Heian period, battles were mainly fought using yumi (longbow) on horseback, but in the late Heian period, battles on foot began to increase and naginata also came to be used on the battlefield. The naginata was appreciated because it was a weapon that could maintain an optimum distance from the enemy in close combat. During the Genpei War (1180–1185), in which the Taira clan was pitted against the Minamoto clan, the naginata rose to a position of particularly high esteem, being regarded as an extremely effective weapon by warriors. The Tale of the Heike, which records the Genpei War, there are descriptions such as ō naginata (lit. big naginata) and ko naginata (lit. little naginata), which show that naginata of various lengths were used. The naginata proved excellent at dismounting cavalry and disabling riders. The widespread adoption of the naginata as a battlefield weapon forced the introduction of greaves as a part of Japanese armor. Ōyamazumi Shrine houses two naginata that are said to have been dedicated by Tomoe Gozen and Benkei at the end of the Heian period and they are designated as Important Cultural Property. However, according to Karl Friday, there were various notations for naginata in the Heian period and the earliest physical evidence for naginata was in the middle of the Kamakura period, so there is a theory that says when they first appeared is unclear. Earlier 10th through 12th century sources refer to "long swords" that while a common medieval term or orthography for naginata, could also simply be referring to conventional swords; one source describes a naginata being drawn with the verb , commonly associated with swords, rather than , the verb otherwise used in medieval texts for unsheathing naginata. Some 11th and 12th century mentions of hoko may actually have been referring to naginata. The commonly assumed association of the naginata and the sōhei is also unclear. Artwork from the late-13th and 14th centuries depict the sōhei with naginata but do not appear to place any special significance to it: the weapons appear as just one of a number of others carried by the monks, and are used by samurai and commoners as well. Depictions of naginata-armed sōhei in earlier periods were created centuries after the fact, and are likely using the naginata as a symbol to distinguish the sōhei from other warriors, rather than giving an accurate portrayal of the events. After the Ōnin War (1467–1477) in the Muromachi period, large-scale group battles started in which mobilized (foot soldiers) fought on foot and in close quarters, and (spear), (longbow), and (Japanese matchlock) became the main weapons. This made and obsolete on the battlefield, and they were often replaced with the and short, lightweight . In the Edo period (1603–1867), the hilts of were often cut off and made into or . This practice of cutting off the hilt of an , , , or and remaking it into a shorter or due to changes in tactics is called and was common in Japan at the time. In Japan there is a saying about swords: "No sword made by modifying a or a is dull in cutting" (薙刀（長巻）直しに鈍刀なし). The meaning of this saying is that and are equipment for actual combat, not works of art or offerings to the , and that the sharpness and durability of swords made from their modifications have been proven on the battlefield. In the peaceful Edo period, weapons' value as battlefield weapons became diminished and their value for martial arts and self-defense rose. The naginata was accepted as a status symbol and self-defense weapon for women of nobility, resulting in the image that "the Naginata is the main weapon used by women". In the Meiji era, it gained popularity along sword martial arts. From the Taisho era to the post-War era, the naginata became popular as a martial art for women, mainly due to the influence of government policies. Although associated with considerably smaller numbers of practitioners, a number of "koryu bujutsu" systems (traditional martial arts) which include older and more combative forms of naginatajutsu remain existent, including Suio Ryu, Araki Ryu, Tendo Ryu, Jikishinkage ryu, Higo Koryu, Tenshin Shoden Katori Shinto Ryu, Toda-ha Buko Ryu, and Yoshin ryu, some of which have authorized representatives outside Japan. Contemporary construction In contemporary naginatajutsu, two types of practice naginata are in common use. The naginata used in atarashii naginata (新しいなぎなた), the shiai-yo, has an oak shaft and a bamboo "blade" (habu). It is used for practice, forms competitions, and sparring. It is between and in length and must weigh over . The "blade" is replaceable. They are often broken or damaged during sparring and can be quickly replaced, being attached to the shaft with tape. The naginata used by koryū practitioners has an oak shaft and blade, carved from a single piece of wood, and may incorporate a disc-shaped guard (tsuba). It is called a kihon-yo. Contemporary usage Naginata can be used to batter, stab, or hook an opponent, but due to their relatively balanced center of mass, are often spun and turned to proscribe a large radius of reach. The curved blade provides a long cutting surface without increasing the overall length of the weapon. Historically, the naginata was often used by foot soldiers to create space on the battlefield. They have several situational advantages over a sword. Their reach is longer, allowing the wielder to keep out of the reach of opponents. The weight of the weapon gave power to strikes and cuts, even though the weight of the weapon is usually thought of as a disadvantage. The weight at the end of the shaft (ishizuki), and the shaft itself (ebu) can be used offensively and defensively. The martial art of wielding the naginata is known as naginatajutsu. Most naginata practice today is in a modernised form, a gendai budō called atarashii Naginata ("new Naginata"), which is organized into regional, national, and international federations, who hold competitions and award ranks. Use of the naginata is also taught within the Bujinkan and in some koryū schools such as Suio Ryu and Tendō-ryū. Naginata practitioners wear an uwagi, obi, and hakama, similar to that worn by kendo practitioners, although the uwagi is generally white. For sparring, armor known as bōgu is worn. Bōgu for naginatajutsu adds and the have a singulated index finger, unlike the mitten-style gloves used for kendo. Gallery	Technology	Polearms	null
153797	https://en.wikipedia.org/wiki/Orion%20%28constellation%29	Orion (constellation)	Orion is a prominent set of stars visible during winter in the northern celestial hemisphere. It is one of the 88 modern constellations; it was among the 48 constellations listed by the 2nd-century astronomer Ptolemy. It is named after a hunter in Greek mythology. Orion is most prominent during winter evenings in the Northern Hemisphere, as are five other constellations that have stars in the Winter Hexagon asterism. Orion's two brightest stars, Rigel (β) and Betelgeuse (α), are both among the brightest stars in the night sky; both are supergiants and slightly variable. There are a further six stars brighter than magnitude 3.0, including three making the short straight line of the Orion's Belt asterism. Orion also hosts the radiant of the annual Orionids, the strongest meteor shower associated with Halley's Comet, and the Orion Nebula, one of the brightest nebulae in the sky. Characteristics Orion is bordered by Taurus to the northwest, Eridanus to the southwest, Lepus to the south, Monoceros to the east, and Gemini to the northeast. Covering 594 square degrees, Orion ranks twenty-sixth of the 88 constellations in size. The constellation boundaries, as set by Belgian astronomer Eugène Delporte in 1930, are defined by a polygon of 26 sides. In the equatorial coordinate system, the right ascension coordinates of these borders lie between and , while the declination coordinates are between and . The constellation's three-letter abbreviation, as adopted by the International Astronomical Union in 1922, is "Ori". Orion is most visible in the evening sky from January to April, winter in the Northern Hemisphere, and summer in the Southern Hemisphere. In the tropics (less than about 8° from the equator), the constellation transits at the zenith. In the period May–July (summer in the Northern Hemisphere, winter in the Southern Hemisphere), Orion is in the daytime sky and thus invisible at most latitudes. However, for much of Antarctica in the Southern Hemisphere's winter months, the Sun is below the horizon even at midday. Stars (and thus Orion, but only the brightest stars) are then visible at twilight for a few hours around local noon, just in the brightest section of the sky low in the North where the Sun is just below the horizon. At the same time of day at the South Pole itself (Amundsen–Scott South Pole Station), Rigel is only 8° above the horizon, and the Belt sweeps just along it. In the Southern Hemisphere's summer months, when Orion is normally visible in the night sky, the constellation is actually not visible in Antarctica because the sun does not set at that time of year south of the Antarctic Circle. In countries close to the equator (e.g., Kenya, Indonesia, Colombia, Ecuador), Orion appears overhead in December around midnight and in the February evening sky. Navigational aid Orion is very useful as an aid to locating other stars. By extending the line of the Belt southeastward, Sirius (α CMa) can be found; northwestward, Aldebaran (α Tau). A line eastward across the two shoulders indicates the direction of Procyon (α CMi). A line from Rigel through Betelgeuse points to Castor and Pollux (α Gem and β Gem). Additionally, Rigel is part of the Winter Circle asterism. Sirius and Procyon, which may be located from Orion by following imaginary lines (see map), also are points in both the Winter Triangle and the Circle. Features Orion's seven brightest stars form a distinctive hourglass-shaped asterism, or pattern, in the night sky. Four stars—Rigel, Betelgeuse, Bellatrix, and Saiph—form a large roughly rectangular shape, at the center of which lies the three stars of Orion's Belt—Alnitak, Alnilam, and Mintaka. His head is marked by an additional 8th star called Meissa, which is fairly bright to the observer. Descending from the "belt" is a smaller line of three stars, Orion's Sword (the middle of which is in fact not a star but the Orion Nebula), also known as the hunter's sword. Many of the stars are luminous hot blue supergiants, with the stars of the belt and sword forming the Orion OB1 association. Standing out by its red hue, Betelgeuse may nevertheless be a runaway member of the same group. Bright stars Betelgeuse, also designated Alpha Orionis, is a massive M-type red supergiant star nearing the end of its life. It is the second brightest star in Orion, and is a semiregular variable star. It serves as the "right shoulder" of the hunter (assuming that he is facing the observer). It is generally the eleventh brightest star in the night sky, but this has varied between being the tenth brightest to the 23rd brightest by the end of 2019. The end of its life is expected to result in a supernova explosion that will be highly visible from Earth, possibly outshining the Earth's moon and being visible during the day. This is most likely to occur within the next 100,000 years. Rigel, also known as Beta Orionis, is a B-type blue supergiant that is the seventh brightest star in the night sky. Similar to Betelgeuse, Rigel is fusing heavy elements in its core and will pass its supergiant stage soon (on an astronomical timescale), either collapsing in the case of a supernova or shedding its outer layers and turning into a white dwarf. It serves as the left foot of the hunter. Bellatrix is designated Gamma Orionis by Johann Bayer. It is the twenty-seventh brightest star in the night sky. Bellatrix is considered a B-type blue giant, though it is too small to explode in a supernova. Bellatrix's luminosity is derived from its high temperature rather than a large radius. Bellatrix marks Orion's left shoulder and it means the "female warrior", and is sometimes known colloquially as the "Amazon Star". It is the closest major star in Orion at only 244.6 light years from our solar system. Mintaka is designated Delta Orionis, despite being the faintest of the three stars in Orion's Belt. Its name means "the belt". It is a multiple star system, composed of a large B-type blue giant and a more massive O-type main-sequence star. The Mintaka system constitutes an eclipsing binary variable star, where the eclipse of one star over the other creates a dip in brightness. Mintaka is the westernmost of the three stars of Orion's Belt, as well as the northernmost. Alnilam is designated Epsilon Orionis and is named for the Arabic phrase meaning "string of pearls". It is the middle and brightest of the three stars of Orion's Belt. Alnilam is a B-type blue supergiant; despite being nearly twice as far from the Sun as the other two belt stars, its luminosity makes it nearly equal in magnitude. Alnilam is losing mass quickly, a consequence of its size. It is the farthest major star in Orion at 1,344 light years. Alnitak, meaning "the girdle", is designated Zeta Orionis, and is the easternmost star in Orion's Belt. It is a triple star system, with the primary star being a hot blue supergiant and the brightest class O star in the night sky. Saiph is designated Kappa Orionis by Bayer, and serves as Orion's right foot. It is of a similar distance and size to Rigel, but appears much fainter. It means the "sword of the giant" Meissa is designated Lambda Orionis, forms Orion's head, and is a multiple star with a combined apparent magnitude of 3.33. Its name means the "shining one". Belt Orion's Belt or The Belt of Orion is an asterism within the constellation. It consists of the three bright stars Zeta (Alnitak), Epsilon (Alnilam), and Delta (Mintaka). Alnitak is around 800 light years away from earth and is 100,000 times more luminous than the Sun and shines with magnitude 1.8; much of its radiation is in the ultraviolet range, which the human eye cannot see. Alnilam is approximately 2,000 light years away from Earth, shines with magnitude 1.70, and with ultraviolet light is 375,000 times more luminous than the Sun. Mintaka is 915 light years away and shines with magnitude 2.21. It is 90,000 times more luminous than the Sun and is a double star: the two orbit each other every 5.73 days. In the Northern Hemisphere, Orion's Belt is best visible in the night sky during the month of January around 9:00 pm, when it is approximately around the local meridian. Just southwest of Alnitak lies Sigma Orionis, a multiple star system composed of five stars that have a combined apparent magnitude of 3.7 and lying 1150 light years distant. Southwest of Mintaka lies the quadruple star Eta Orionis. Sword Orion's Sword contains the Orion Nebula, the Messier 43 nebula, the Running Man Nebula, and the stars Theta Orionis, Iota Orionis, and 42 Orionis. Head Three stars comprise a small triangle that marks the head. The apex is marked by Meissa (Lambda Orionis), a hot blue giant of spectral type O8 III and apparent magnitude 3.54, which lies some 1100 light years distant. Phi-1 and Phi-2 Orionis make up the base. Also nearby is the very young star FU Orionis. Club Stretching north from Betelgeuse are the stars that make up Orion's club. Mu Orionis marks the elbow, Nu and Xi mark the handle of the club, and Chi1 and Chi2 mark the end of the club. Just east of Chi1 is the Mira-type variable red giant U Orionis. Shield West from Bellatrix lie six stars all designated Pi Orionis (π1 Ori, π2 Ori, π3 Ori, π4 Ori, π5 Ori and π6 Ori) which make up Orion's shield. Meteor showers Around 20 October each year the Orionid meteor shower (Orionids) reaches its peak. Coming from the border with the constellation Gemini as many as 20 meteors per hour can be seen. The shower's parent body is Halley's Comet. Deep-sky objects Hanging from Orion's belt is his sword, consisting of the multiple stars θ1 and θ2 Orionis, called the Trapezium and the Orion Nebula (M42). This is a spectacular object that can be clearly identified with the naked eye as something other than a star. Using binoculars, its clouds of nascent stars, luminous gas, and dust can be observed. The Trapezium cluster has many newborn stars, including several brown dwarfs, all of which are at an approximate distance of 1,500 light-years. Named for the four bright stars that form a trapezoid, it is largely illuminated by the brightest stars, which are only a few hundred thousand years old. Observations by the Chandra X-ray Observatory show both the extreme temperatures of the main stars—up to 60,000 kelvins—and the star forming regions still extant in the surrounding nebula. M78 (NGC 2068) is a nebula in Orion. With an overall magnitude of 8.0, it is significantly dimmer than the Great Orion Nebula that lies to its south; however, it is at approximately the same distance, at 1600 light-years from Earth. It can easily be mistaken for a comet in the eyepiece of a telescope. M78 is associated with the variable star V351 Orionis, whose magnitude changes are visible in very short periods of time. Another fairly bright nebula in Orion is NGC 1999, also close to the Great Orion Nebula. It has an integrated magnitude of 10.5 and is 1500 light-years from Earth. The variable star V380 Orionis is embedded in NGC 1999. Another famous nebula is IC 434, the Horsehead Nebula, near ζ Orionis. It contains a dark dust cloud whose shape gives the nebula its name. NGC 2174 is an emission nebula located 6400 light-years from Earth. Besides these nebulae, surveying Orion with a small telescope will reveal a wealth of interesting deep-sky objects, including M43, M78, as well as multiple stars including Iota Orionis and Sigma Orionis. A larger telescope may reveal objects such as the Flame Nebula (NGC 2024), as well as fainter and tighter multiple stars and nebulae. Barnard's Loop can be seen on very dark nights or using long-exposure photography. All of these nebulae are part of the larger Orion molecular cloud complex, which is located approximately 1,500 light-years away and is hundreds of light-years across. It is one of the most intense regions of stellar formation visible within our galaxy. History and mythology The distinctive pattern of Orion is recognized in numerous cultures around the world, and many myths are associated with it. Orion is used as a symbol in the modern world. Ancient Near East The Babylonian star catalogues of the Late Bronze Age name Orion , "The Heavenly Shepherd" or "True Shepherd of Anu" – Anu being the chief god of the heavenly realms. The Babylonian constellation is sacred to Papshukal and Ninshubur, both minor gods fulfilling the role of 'messenger to the gods'. Papshukal is closely associated with the figure of a walking bird on Babylonian boundary stones, and on the star map the figure of the Rooster is located below and behind the figure of the True Shepherd—both constellations represent the herald of the gods, in his bird and human forms respectively. In ancient Egypt, the stars of Orion were regarded as a god, called Sah. Because Orion rises before Sirius, the star whose heliacal rising was the basis for the Solar Egyptian calendar, Sah was closely linked with Sopdet, the goddess who personified Sirius. The god Sopdu is said to be the son of Sah and Sopdet. Sah is syncretized with Osiris, while Sopdet is syncretized with Osiris' mythological wife, Isis. In the Pyramid Texts, from the 24th and 23rd centuries BC, Sah is one of many gods whose form the dead pharaoh is said to take in the afterlife. The Armenians identified their legendary patriarch and founder Hayk with Orion. Hayk is also the name of the Orion constellation in the Armenian translation of the Bible. The Bible mentions Orion three times, naming it "Kesil" (כסיל, literally – fool). Though, this name perhaps is etymologically connected with "Kislev", the name for the ninth month of the Hebrew calendar (i.e. November–December), which, in turn, may derive from the Hebrew root K-S-L as in the words "kesel, kisla" (כֵּסֶל, כִּסְלָה, hope, positiveness), i.e. hope for winter rains.: Job 9:9 ("He is the maker of the Bear and Orion"), Job 38:31 ("Can you loosen Orion's belt?"), and Amos 5:8 ("He who made the Pleiades and Orion"). In ancient Aram, the constellation was known as Nephîlā′, the Nephilim are said to be Orion's descendants. Greco-Roman antiquity In Greek mythology, Orion was a gigantic, supernaturally strong hunter, born to Euryale, a Gorgon, and Poseidon (Neptune), god of the sea. One myth recounts Gaia's rage at Orion, who dared to say that he would kill every animal on Earth. The angry goddess tried to dispatch Orion with a scorpion. This is given as the reason that the constellations of Scorpius and Orion are never in the sky at the same time. However, Ophiuchus, the Serpent Bearer, revived Orion with an antidote. This is said to be the reason that the constellation of Ophiuchus stands midway between the Scorpion and the Hunter in the sky. The constellation is mentioned in Horace's Odes (Ode 3.27.18), Homer's Odyssey (Book 5, line 283) and Iliad, and Virgil's Aeneid (Book 1, line 535) Middle East In medieval Muslim astronomy, Orion was known as al-jabbar, "the giant". Orion's sixth brightest star, Saiph, is named from the Arabic, saif al-jabbar, meaning "sword of the giant". China In China, Orion was one of the 28 lunar mansions Sieu (Xiù) (宿). It is known as Shen (參), literally meaning "three", for the stars of Orion's Belt. (See Chinese constellations) The Chinese character 參 (pinyin shēn) originally meant the constellation Orion (); its Shang dynasty version, over three millennia old, contains at the top a representation of the three stars of Orion's belt atop a man's head (the bottom portion representing the sound of the word was added later). India The Rigveda refers to the Orion Constellation as Mriga (The Deer). Nataraja, 'the cosmic dancer', is often interpreted as the representation of Orion. Rudra, the Rigvedic form of Shiva, is the presiding deity of Ardra nakshatra (Betelgeuse) of Hindu astrology. The Jain Symbol carved in Udayagiri and Khandagiri Caves, India in 1st century BCE has striking resemblance with Orion. Bugis sailors identified the three stars in Orion's Belt as tanra tellué, meaning "sign of three". European folklore In old Hungarian tradition, Orion is known as "Archer" (Íjász), or "Reaper" (Kaszás). In recently rediscovered myths, he is called Nimrod (Hungarian: Nimród), the greatest hunter, father of the twins Hunor and Magor. The π and o stars (on upper right) form together the reflex bow or the lifted scythe. In other Hungarian traditions, Orion's belt is known as "Judge's stick" (Bírópálca). In Scandinavian tradition, Orion's belt was known as "Frigg's Distaff" (friggerock) or "Freyja's distaff". The Finns call Orion's belt and the stars below it "Väinämöinen's scythe" (Väinämöisen viikate). Another name for the asterism of Alnilam, Alnitak and Mintaka is "Väinämöinen's Belt" (Väinämöisen vyö) and the stars "hanging" from the belt as "Kaleva's sword" (Kalevanmiekka). In Siberia, the Chukchi people see Orion as a hunter; an arrow he has shot is represented by Aldebaran (Alpha Tauri), with the same figure as other Western depictions. There are claims in popular media that the Adorant from the Geißenklösterle cave, an ivory carving estimated to be 35,000 to 40,000 years old, is the first known depiction of the constellation. Scholars dismiss such interpretations, saying that perceived details such as a belt and sword derive from preexisting features in the grain structure of the ivory. Americas The Seri people of northwestern Mexico call the three stars in the belt of Orion Hapj (a name denoting a hunter) which consists of three stars: Hap (mule deer), Haamoja (pronghorn), and Mojet (bighorn sheep). Hap is in the middle and has been shot by the hunter; its blood has dripped onto Tiburón Island. The same three stars are known in Spain and most of Latin America as "Las tres Marías" (Spanish for "The Three Marys"). In Puerto Rico, the three stars are known as the "Los Tres Reyes Magos" (Spanish for The three Wise Men). The Ojibwa (Chippewa) Native Americans call this constellation Kabibona'kan, the Winter Maker, as its presence in the night sky heralds winter. To the Lakota Native Americans, Tayamnicankhu (Orion's Belt) is the spine of a bison. The great rectangle of Orion is the bison's ribs; the Pleiades star cluster in nearby Taurus is the bison's head; and Sirius in Canis Major, known as Tayamnisinte, is its tail. Another Lakota myth mentions that the bottom half of Orion, the Constellation of the Hand, represented the arm of a chief that was ripped off by the Thunder People as a punishment from the gods for his selfishness. His daughter offered to marry the person who can retrieve his arm from the sky, so the young warrior Fallen Star (whose father was a star and whose mother was human) returned his arm and married his daughter, symbolizing harmony between the gods and humanity with the help of the younger generation. The index finger is represented by Rigel; the Orion Nebula is the thumb; the Belt of Orion is the wrist; and the star Beta Eridani is the pinky finger. Austronesian The seven primary stars of Orion make up the Polynesian constellation Heiheionakeiki which represents a child's string figure similar to a cat's cradle. Several precolonial Filipinos referred to the belt region in particular as "balatik" (ballista) as it resembles a trap of the same name which fires arrows by itself and is usually used for catching pigs from the bush. Spanish colonization later led to some ethnic groups referring to Orion's belt as "Tres Marias" or "Tatlong Maria." In Māori tradition, the star Rigel (known as Puanga or Puaka) is closely connected with the celebration of Matariki. The rising of Matariki (the Pleiades) and Rigel before sunrise in midwinter marks the start of the Māori year. In Javanese culture, the constellation is often called Lintang Waluku or Bintang Bajak, referring to the shape of a paddy field plow. Contemporary symbolism The imagery of the belt and sword has found its way into popular western culture, for example in the form of the shoulder insignia of the 27th Infantry Division of the United States Army during both World Wars, probably owing to a pun on the name of the division's first commander, Major General John F. O'Ryan. The film distribution company Orion Pictures used the constellation as its logo. Depictions In artistic renderings, the surrounding constellations are sometimes related to Orion: he is depicted standing next to the river Eridanus with his two hunting dogs Canis Major and Canis Minor, fighting Taurus. He is sometimes depicted hunting Lepus the hare. He sometimes is depicted to have a lion's hide in his hand. There are alternative ways to visualise Orion. From the Southern Hemisphere, Orion is oriented south-upward, and the belt and sword are sometimes called the saucepan or pot in Australia and New Zealand. Orion's Belt is called Drie Konings (Three Kings) or the Drie Susters (Three Sisters) by Afrikaans speakers in South Africa and are referred to as les Trois Rois (the Three Kings) in Daudet's Lettres de Mon Moulin (1866). The appellation Driekoningen (the Three Kings) is also often found in 17th- and 18th-century Dutch star charts and seaman's guides. The same three stars are known in Spain, Latin America, and the Philippines as "Las Tres Marías" (The Three Marys), and as "Los Tres Reyes Magos" (The three Wise Men) in Puerto Rico. Even traditional depictions of Orion have varied greatly. Cicero drew Orion in a similar fashion to the modern depiction. The Hunter held an unidentified animal skin aloft in his right hand; his hand was represented by Omicron2 Orionis and the skin was represented by the 5 stars designated Pi Orionis. Kappa and Beta Orionis represented his left and right knees, while Eta and Lambda Leporis were his left and right feet, respectively. As in the modern depiction, Delta, Epsilon, and Zeta represented his belt. His left shoulder was represented by Alpha Orionis, and Mu Orionis made up his left arm. Lambda Orionis was his head and Gamma, his right shoulder. The depiction of Hyginus was similar to that of Cicero, though the two differed in a few important areas. Cicero's animal skin became Hyginus's shield (Omicron and Pi Orionis), and instead of an arm marked out by Mu Orionis, he holds a club (Chi Orionis). His right leg is represented by Theta Orionis and his left leg is represented by Lambda, Mu, and Epsilon Leporis. Further Western European and Arabic depictions have followed these two models. Future Orion is located on the celestial equator, but it will not always be so located due to the effects of precession of the Earth's axis. Orion lies well south of the ecliptic, and it only happens to lie on the celestial equator because the point on the ecliptic that corresponds to the June solstice is close to the border of Gemini and Taurus, to the north of Orion. Precession will eventually carry Orion further south, and by AD 14000, Orion will be far enough south that it will no longer be visible from the latitude of Great Britain. Further in the future, Orion's stars will gradually move away from the constellation due to proper motion. However, Orion's brightest stars all lie at a large distance from the Earth on an astronomical scale—much farther away than Sirius, for example. Orion will still be recognizable long after most of the other constellations—composed of relatively nearby stars—have distorted into new configurations, with the exception of a few of its stars eventually exploding as supernovae, for example Betelgeuse, which is predicted to explode sometime in the next million years.	Physical sciences	Constellations	null
153859	https://en.wikipedia.org/wiki/Palisade	Palisade	A palisade, sometimes called a stakewall or a paling, is typically a row of closely placed, high vertical standing tree trunks or wooden or iron stakes used as a fence for enclosure or as a defensive wall. Palisades can form a stockade. Etymology Palisade derives from pale, from the Latin word , meaning stake, specifically when used side by side to create a wood defensive wall. (see 'pale', English: Etymology 2 on Wiktionary). Typical construction Typical construction consisted of small or mid-sized tree trunks aligned vertically, with as little free space in between as possible. The trunks were sharpened or pointed at the top, and were driven into the ground and sometimes reinforced with additional construction. The height of a palisade ranged from around a metre to as high as 3–4 m. As a defensive structure, palisades were often used in conjunction with earthworks. Palisades were an excellent option for small forts or other hastily constructed fortifications. Since they were made of wood, they could often be quickly and easily built from readily available materials. They proved to be effective protection for short-term conflicts and were an effective deterrent against small forces. However, because they were wooden constructions they were also vulnerable to fire and siege weapons. Often, a palisade would be constructed around a castle as a temporary wall until a permanent stone wall could be erected. Ancient Greece and Rome Both the Greeks and Romans created palisades to protect their military camps. The Roman historian Livy describes the Greek method as being inferior to that of the Romans during the Second Macedonian War. The Greek stakes were too large to be easily carried and were spaced too far apart. This made it easy for enemies to uproot them and create a large enough gap in which to enter. In contrast, the Romans used smaller and easier to carry stakes which were placed closer together, making them more difficult to uproot. Precolumbian North America The Iroquoian peoples, who coalesced as tribes around the Great Lakes, often defended their settlements with palisades. Within the palisades the peoples lived in communal groups in numerous longhouses, sometimes in communities as large as 2,000 people. Archeological evidence of such palisades has been found at numerous 15th and 16th-century sites in both Ontario, Canada, and in New York, United States. Many settlements of the native Mississippian culture of the Midwestern United States used palisades. A prominent example is the Cahokia Mounds site in Collinsville, Illinois. A wooden stockade with a series of watchtowers or bastions at regular intervals formed a enclosure around Monk's Mound and the Grand Plaza. Archaeologists found evidence of the stockade during excavation of the area and indications that it was rebuilt several times, in slightly different locations. The stockade seems to have separated Cahokia's main ceremonial precinct from other parts of the city, as well as being a defensive structure. Other examples include the Angel Mounds site in southern Indiana, Aztalan State Park in Wisconsin, the Kincaid site in Illinois, the Parkin site and the Nodena sites in northeastern Arkansas, and the Etowah site in Georgia. Colonial America Palisaded settlements were common in Colonial North America, for protection against indigenous peoples and wild animals. The English settlements in Jamestown, Virginia (1607), Cupids, Newfoundland (1610) and Plymouth, Massachusetts (1620) were all originally fortifications that were surrounded by palisades. Such defensive palisades were also frequently used in New France. In addition, colonial architecture used vertical palings as the walls of houses, in what was called poteaux en terre construction. Some 18th-century houses in this style survive in Ste. Genevieve, Missouri, initially settled by French colonists from the Illinois Country to the east of the Mississippi River. Ottoman Empire A "palanka" was a type of wooden fort constructed of palisades, built by the Ottoman Empire in the Balkans during the 16th and 17th centuries. They could be erected for a variety of reasons such as protecting a strategically valuable area or a town Some palankas evolved into larger settlements. Half-timber palisade construction In the late nineteenth century, when milled lumber was not available or practical, many Adirondack buildings were built using a palisade architecture. The walls were made of vertical half timbers; the outside, rounded half with its bark still on faced Adirondack weather, while the inside half was sanded and varnished for a finished wood look. Typically, the cracks between the vertical logs were filled with moss and sometimes covered with small sticks. Inside, the cracks were covered with narrow wooden battens. This palisade style was much more efficient to build than the traditional horizontal log cabin, since two half logs provided more surface area than one whole log and the vertical alignment meant a stronger structure for supporting loads like upper stories and roofs. It also presented a more finished look inside. Examples of this architectural style can still be found in the Adirondacks, such as around Big Moose Lake. Modern uses In areas with extremely high rates of violent crime and property theft, a common means to prevent crime is for residential houses to be protected by perimeter defenses such as ornamental iron bars, brick walls, steel palisade fences, wooden palisade fences and electrified palisade fences (railings). The City of Johannesburg promotes the use of palisade fencing over opaque, usually brick, walls, as criminals cannot hide as easily behind the fence. Its manual on safety includes guidance, such as avoiding having vegetation alongside the fence, as this allows criminals to make an unseen breach.	Technology	Fortification	null
153861	https://en.wikipedia.org/wiki/Moat	Moat	A moat is a deep, broad ditch dug around a castle, fortification, building, or town, historically to provide it with a preliminary line of defence. Moats can be dry or filled with water. In some places, moats evolved into more extensive water defences, including natural or artificial lakes, dams and sluices. In older fortifications, such as hillforts, they are usually referred to simply as ditches, although the function is similar. In later periods, moats or water defences may be largely ornamental. They could also act as a sewer. Historical use Ancient Some of the earliest evidence of moats has been uncovered around ancient Egyptian fortresses. One example is at Buhen, a settlement excavated in Nubia. Other evidence of ancient moats is found in the ruins of Babylon, and in reliefs from ancient Egypt, Assyria, and other cultures in the region. Evidence of early moats around settlements has been discovered in many archaeological sites throughout Southeast Asia, including Noen U-Loke, Ban Non Khrua Chut, Ban Makham Thae and Ban Non Wat. The use of the moats could have been either for defensive or agriculture purposes. Medieval Moats were excavated around castles and other fortifications as part of the defensive system as an obstacle immediately outside the walls. In suitable locations, they might be filled with water. A moat made access to the walls difficult for siege weapons such as siege towers and battering rams, which needed to be brought up against a wall to be effective. A water-filled moat made the practice of mining – digging tunnels under the castles in order to effect a collapse of the defences – very difficult as well. Segmented moats have one dry section and one section filled with water. Dry moats that cut across the narrow part of a spur or peninsula are called neck ditches. Moats separating different elements of a castle, such as the inner and outer wards, are cross ditches. The word was adapted in Middle English from the Old French () and was first applied to the central mound on which a castle was erected (see Motte and bailey) and then came to be applied to the excavated ring, a 'dry moat'. The shared derivation implies that the two features were closely related and possibly constructed at the same time. The term moat is also applied to natural formations reminiscent of the artificial structure and to similar modern architectural features. Later western fortification With the introduction of siege artillery, a new style of fortification emerged in the 16th century using low walls and projecting strong points called bastions, which was known as the trace italienne. The walls were further protected from infantry attack by wet or dry moats, sometimes in elaborate systems. When this style of fortification was superseded by lines of polygonal forts in the mid-19th century, moats continued to be used for close protection. Africa The Walls of Benin were a combination of ramparts and moats, called Iya, used as a defence of the capital Benin City in present-day Edo State of Nigeria. It was considered the largest man-made structure lengthwise, second only to the Great Wall of China and the largest earthwork in the world. Recent work by Patrick Darling has established it as the largest man-made structure in the world, larger than Sungbo's Eredo, also in Nigeria. It enclosed of community lands. Its length was over of earth boundaries. It was estimated that earliest construction began in 800 and continued into the mid-15th century. The walls are built of a ditch and dike structure, the ditch dug to form an inner moat with the excavated earth used to form the exterior rampart. The Benin Walls were ravaged by the British in 1897. Scattered pieces of the walls remain in Edo, with material being used by the locals for building purposes. The walls continue to be torn down for real-estate developments. The Walls of Benin City were the world's largest man-made structure. Fred Pearce wrote in New Scientist:They extend for some 16,000 kilometres in all, in a mosaic of more than 500 interconnected settlement boundaries. They cover 6,500 square kilometres and were all dug by the Edo people. In all, they are four times longer than the Great Wall of China, and consumed a hundred times more material than the Great Pyramid of Cheops. They took an estimated 150 million hours of digging to construct, and are perhaps the largest single archaeological phenomenon on the planet. Asia Japanese castles often have very elaborate moats, with up to three moats laid out in concentric circles around the castle and a host of different patterns engineered around the landscape. The outer moat of a Japanese castle typically protects other support buildings in addition to the castle. As many Japanese castles have historically been a very central part of their cities, the moats have provided a vital waterway to the city. Even in modern times the moat system of the Tokyo Imperial Palace consists of a very active body of water, hosting everything from rental boats and fishing ponds to restaurants. Most modern Japanese castles have moats filled with water, but castles in the feudal period more commonly had 'dry moats' , a trench. A is a dry moat dug into a slope. A is a series of parallel trenches running up the sides of the excavated mountain, and the earthen wall, which was also called , was an outer wall made of earth dug out from a moat. Even today it is common for mountain Japanese castles to have dry moats. A is a moat filled with water. Moats were also used in the Forbidden City and Xi'an in China; in Vellore Fort in India; Hsinchu in Taiwan; and in Southeast Asia, such as at Angkor Wat in Cambodia; Mandalay in Myanmar; Chiang Mai in Thailand and Huế in Vietnam. Australia The only moated fort ever built in Australia was Fort Lytton in Brisbane. As Brisbane was much more vulnerable to attack than either Sydney or Melbourne a series of coastal defences was built throughout Moreton Bay, Fort Lytton being the largest. Built between 1880 and 1881 in response to fear of a Russian invasion, it is a pentagonal fortress concealed behind grassy embankments and surrounded by a water-filled moat. North America Moats were developed independently by North American indigenous people of the Mississippian culture as the outer defence of some fortified villages. The remains of a 16th-century moat are still visible at the Parkin Archeological State Park in eastern Arkansas. The Maya people also used moats, for example in the city of Becan. European colonists in the Americas often built dry ditches surrounding forts built to protect important landmarks, harbours or cities (e.g. Fort Jay on Governors Island in New York Harbor). Photo gallery Modern usage Architectural usage Dry moats were a key element used in French Classicism and Beaux-Arts architecture dwellings, both as decorative designs and to provide discreet access for service. Excellent examples of these can be found in Newport, Rhode Island at Miramar (mansion) and The Elms, as well as at Carolands, outside of San Francisco, California, and at Union Station in Toronto, Ontario, Canada. Additionally, a dry moat can allow light and fresh air to reach basement workspaces, as for example at the James Farley Post Office in New York City. Anti-terrorist moats Whilst moats are no longer a significant tool of warfare, modern architectural building design continues to use them as a defence against certain modern threats, such as terrorist attacks from car bombs and improvised fighting vehicles. For example, the new location of the Embassy of the United States in London, opened in 2018, includes a moat among its security features - the first moat built in England for more than a century. Modern moats may also be used for aesthetic or ergonomic purposes. The Catawba Nuclear Station has a concrete moat around the sides of the plant not bordering a lake. The moat is a part of precautions added to such sites after the September 11, 2001 attacks. Safety moats Moats, rather than fences, separate animals from spectators in many modern zoo installations. Moats were first used in this way by Carl Hagenbeck at his Tierpark in Hamburg, Germany. The structure, with a vertical outer retaining wall rising direct from the moat, is an extended usage of the ha-ha of English landscape gardening. Border defence moats In 2004, plans were suggested for a two-mile moat across the southern border of the Gaza Strip to prevent tunnelling from Egyptian territory to the border town of Rafah. In 2008, city officials in Yuma, Arizona planned to dig out a two-mile stretch of a wetland known as Hunters Hole to control immigrants coming from Mexico. Pest control moats Researchers of jumping spiders, which have excellent vision and adaptable tactics, built water-filled miniature moats, too wide for the spiders to jump across. Some specimens were rewarded for jumping then swimming and others for swimming only. Portia fimbriata from Queensland generally succeeded, for whichever method they were rewarded. When specimens from two different populations of Portia labiata were set the same task, members of one population determined which method earned them a reward, whilst members of the other continued to use whichever method they tried first and did not try to adapt. As a basic method of pest control in bonsai, a moat may be used to restrict access of crawling insects to the bonsai.	Technology	Fortification	null
153882	https://en.wikipedia.org/wiki/Spandex	Spandex	Spandex, Lycra, or elastane is a synthetic fiber known for its exceptional elasticity. It is a polyether-polyurea copolymer that was invented in 1958 by chemist Joseph Shivers at DuPont. Name The name spandex, which is an anagram of the word "expands", is the preferred name in North America. In continental Europe, it is referred to by variants of elastane. It is primarily known as Lycra in the UK, Ireland, Portugal, Spain, Latin America, Australia, and New Zealand. Brand names for spandex include Lycra (made by The Lycra Company, previously a division of DuPont Textiles and Interiors), Elaspan (The Lycra Company), Acepora (Taekwang Group), Creora (Hyosung), INVIYA (Indorama Corporation), ROICA and Dorlastan (Asahi Kasei), Linel (Fillattice), and ESPA (Toyobo). Production Unlike many other synthetic fibers, spandex cannot be melt-processed because the polymer degrades upon melting. Spandex fibers are produced by several spinning technologies. Typically, a concentrated solution of the polymer is drawn through spinnerets at temperatures where the solvent evaporates. Spandex is mainly composed of a polyurea derived from the reaction of a diol and a diisocyanate. Two classes of spandex are defined by the "macrodiols". One class of macrodiols is the oligomer produced from tetrahydrofuran (i.e. polytetrahydrofuran). Another class of diols, the so-called ester diols, are oligomers derived from condensation of adipic acid and glycols. Spandex produced from the ester diols is more resilient photochemically and to chlorinated water. Almost always, the diisocyanate is methylenebis(phenyl isocyanate). The key linking reaction is formation or the urea (aka urethane): The polyurea is usually treated with various diamines, which function as chain extenders. Function The exceptional elasticity of spandex fibers increases the clothing's pressure comfort, enhancing the ease of body movements. Pressure comfort is the response towards clothing by the human body's pressure receptors (mechanoreceptors present in skin sensory cells). The sensation response is affected mainly by the stretch, snug, loose, heavy, lightweight, soft, and stiff structure of the material. The elasticity and strength (stretching up to five times its length) of spandex has been incorporated into a wide range of garments, especially in skin-tight garments. A benefit of spandex is its significant strength and elasticity and its ability to return to the original shape after stretching and faster drying than ordinary fabrics. For clothing, spandex is usually mixed with cotton or polyester, and accounts for a small percentage of the final fabric, which therefore retains most of the look and feel of the other fibers. An estimated 80% of clothing sold in the United States contained spandex in 2010. Gallery History The easy condensation of diols and diisocyanates was recognized in the 1930s as the result of work by Otto Bayer. Fibers suitable for replacing nylon were not created from urethanes, but instead this theme led to a family of specialized elastic fabrics. In the post-World War II era, DuPont Textiles Fibers Department, formed in 1952, became the most profitable division of DuPont, dominating the synthetic fiber market worldwide. At this time, women began to emerge as a significant group of consumers because of their need for underwear and hosiery. After conducting market research to find out what women wanted from textiles, DuPont began developing fibers to meet such needs—including a better fiber for women's girdles, which were commonly made of rubber at the time. In the early 1950s chemist Joseph C. Shivers modified Dacron polyester, producing an elastic fiber that could withstand high temperatures. Lycra brand To distinguish its brand of spandex fiber, DuPont chose the trade name Lycra (originally called Fiber K). DuPont launched an extensive publicity campaign for its Lycra brand, taking advertisements and full-page ads in top women's magazines. Audrey Hepburn helped catapult the brand on and off-screen during this time; models and actresses like Joan Collins and Ann-Margret followed Hepburn's aesthetic by posing in Lycra clothing for photo shoots and magazine covers. By the mid-1970s, with the emergence of the women's liberation movement, girdle sales began to drop as they came to be associated with anti-independence and emblematic of an era that was quickly passing away. In response, DuPont marketed Lycra as the aerobic fitness movement emerged in the 1970s. The association of Lycra with fitness had been established at the 1968 Winter Olympic Games, when the French ski team wore Lycra garments. The fiber came to be especially popular in mid-thigh-length shorts worn by cyclists. By the 1980s, the fitness trend had reached its height in popularity and fashionistas began wearing shorts on the street. Spandex proved such a popular fiber in the garment industry that, by 1987, DuPont had trouble meeting worldwide demand. In the 1990s a variety of other items made with spandex proved popular, including a successful line of body-shaping foundation garments sold under the trade name Bodyslimmers. As the decade progressed, shirts, pants, dresses, and even shoes were being made with spandex blends, and mass-market retailers like Banana Republic were even using it for menswear. In 2019, control of the Lycra Company was sold by Koch Industries to Shandong Ruyi. Environmental impact Most clothes containing spandex are difficult to recycle. Even a 5% inclusion of spandex will render the fabric incompatible with most mechanical recycling machines.	Technology	Fabrics and fibers	null
153911	https://en.wikipedia.org/wiki/Invisibility	Invisibility	Invisibility is the state of an object that cannot be seen. An object in this state is said to be invisible (literally, "not visible"). The phenomenon is studied by physics and perceptual psychology. Since objects can be seen by light from a source reflecting off their surfaces and hitting the viewer's eyes, the most natural form of invisibility (whether real or fictional) is an object that neither reflects nor absorbs light (that is, it allows light to pass through it). This is known as transparency, and is seen in many naturally occurring materials (although no naturally occurring material is 100% transparent). Invisibility perception depends on several optical and visual factors. For example, invisibility depends on the eyes of the observer and/or the instruments used. Thus an object can be classified as "invisible" to a person, animal, instrument, etc. In research on sensorial perception it has been shown that invisibility is perceived in cycles. Invisibility is often considered to be the supreme form of camouflage, as it does not reveal to the viewer any kind of vital signs, visual effects, or any frequencies of the electromagnetic spectrum detectable to the human eye, instead making use of radio, infrared or ultraviolet wavelengths. In illusion optics, invisibility is a special case of illusion effects: the illusion of free space. The term is often used in fantasy and science fiction, where objects cannot be seen by means of magic or hypothetical technology. Practical efforts Technology can be used theoretically or practically to render real-world objects invisible. Making use of a real-time image displayed on a wearable display, it is possible to create a see-through effect. This is known as active camouflage. Though stealth technology is declared to be invisible to radar, all officially disclosed applications of the technology can only reduce the size and/or clarity of the signature detected by radar. In 2003 the Chilean scientist Gunther Uhlmann postulates the first mathematical equations to create invisible materials. In 2006, a team effort of researchers from Britain and the US announced the development of a real cloak of invisibility, an artificially made meta material that is invisible to the microwave spectrum, though it is only in its first stages. In filmmaking, people, objects, or backgrounds can be made to look invisible on camera through a process known as chroma keying. Engineers and scientists have performed various kinds of research to investigate the possibility of finding ways to create real optical invisibility (cloaks) for objects. Methods are typically based on implementing the theoretical techniques of transformation optics, which have given rise to several theories of cloaking. Currently, a practical cloaking device does not exist. A 2006 theoretical work predicts that the imperfections are minor, and metamaterials may make real-life "cloaking devices" practical. The technique is predicted to be applied to radio waves within five years, and the distortion of visible light is an eventual possibility. The theory that light waves can be acted upon the same way as radio waves is now a popular idea among scientists. The agent can be compared to a stone in a river, around which water passes, but slightly down-stream leaves no trace of the stone. Comparing light waves to the water, and whatever object that is being "cloaked" to the stone, the goal is to have light waves pass around that object, leaving no visible aspects of it, possibly not even a shadow. This is the technique depicted in the 2000 television portrayal of The Invisible Man. Two teams of scientists worked separately to create two "Invisibility Cloaks" from 'metamaterials' engineered at the nanoscale level. They demonstrated for the first time the possibility of cloaking three-dimensional (3-D) objects with artificially engineered materials that redirect radar, light or other waves around an object. While one uses a type of fishnet of metal layers to reverse the direction of light, the other uses tiny silver wires. Xiang Zhang, of the University of California, Berkeley said: "In the case of invisibility cloaks or shields, the material would need to curve light waves completely around the object like a river flowing around a rock. An observer looking at the cloaked object would then see light from behind it, making it seem to disappear." UC Berkeley researcher Jason Valentine's team made a material that affects light near the visible spectrum, in a region used in fibre optics: 'Instead of the fish appearing to be slightly ahead of where it is in the water, it would actually appear to be above the water's surface. For a metamaterial to produce negative refraction, it must have a structural array smaller than the wavelength of the electromagnetic radiation being used." Valentine's team created their 'fishnet' material by stacking silver and metal dielectric layers on top of each other and then punching holes through them. The other team used an oxide template and grew silver nanowires inside porous aluminum oxide at tiny distances apart, smaller than the wavelength of visible light. This material refracts visible light. The Imperial College London research team achieved results with microwaves. An invisibility cloak layout of a copper cylinder was produced in May, 2008, by physicist Professor Sir John Pendry. Scientists working with him at Duke University in the US put the idea into practice. Pendry, who theorized the invisibility cloak "as a joke" to illustrate the potential of metamaterials, said in an interview in August 2011 that grand, theatrical manifestations of his idea are probably overblown: "I think it’s pretty sure that any cloak that Harry Potter would recognize is not on the table. You could dream up some theory, but the very practicality of making it would be so impossible. But can you hide things from light? Yes. Can you hide things which are a few centimeters across? Yes. Is the cloak really flexible and flappy? No. Will it ever be? No. So you can do quite a lot of things, but there are limitations. There are going to be some disappointed kids around, but there might be a few people in industry who are very grateful for it." In Turkey in 2009, Bilkent University Search Center Of Nanotechnology researches explained and published in New Journal of Physics that they achieved to make invisibility real in practice using nanotechnology making an object invisible with no shadows etc. next to perfect transparent scene by producing nanotechnologic material that can also be produced like a suit anyone can wear. In 2019, Hyperstealth Biotechnology has patented the technology behind a material that bends light to make people and objects near invisible to the naked eye. The material, called Quantum Stealth, is currently still in the prototyping stage, but was developed by the company's CEO Guy Cramer primarily for military purposes, to conceal agents and equipment such as tanks and jets in the field. Unlike traditional camouflage materials, which are limited to specific conditions such as forests or deserts, according to Cramer this "invisibility cloak" works in any environment or season, at any time of day. This is despite its actual application requiring artificial backgrounds made up of horizontal lines. Psychological A person can be described as invisible if others refuse to see them or routinely overlook them. The term was used in this manner in the title of the book Invisible Man, by Ralph Ellison, in reference to the protagonist, likely modeled after the author, being overlooked on account of his status as an African American. This is supported by the quote taken from the Prologue, "I am invisible, understand, simply because people refuse to see me." (Prologue.1) Fictional use In fiction, people or objects can be rendered completely invisible by several means: Magical objects such as rings, cloaks and amulets can be worn to grant the wearer permanent invisibility (or temporary invisibility until the object is taken off). Magical potions can be consumed to grant temporary or permanent invisibility. Magic spells can be cast on people or objects, usually giving temporary invisibility. Some mythical creatures can make themselves invisible at will, such as in some tales in which leprechauns or Chinese dragons can shrink so much that humans cannot see them. In science fiction, the idea of a "cloaking device". In some works, the power of magic creates an effective means of invisibility by distracting anyone who might notice the character. But since the character is not truly invisible, the effect could be betrayed by mirrors or other reflective surfaces. Where magical invisibility is concerned, the issue may arise of whether the clothing worn by and any items carried by the invisible being are also rendered invisible. In general they are also regarded as being invisible, but in some instances clothing remains visible and must be removed for the full invisibility effect.	Physical sciences	Optics	Physics
11688148	https://en.wikipedia.org/wiki/Chlamydia%20%28genus%29	Chlamydia (genus)	Chlamydia is a genus of pathogenic Gram-negative bacteria that are obligate intracellular parasites. Chlamydia infections are the most common bacterial sexually transmitted diseases in humans and are the leading cause of infectious blindness worldwide. Species include Chlamydia trachomatis (a human pathogen), Ch. suis (affects only swine), and Ch. muridarum (affects only mice and hamsters). Humans mainly contract Ch. trachomatis, Ch. pneumoniae, Ch. abortus, and Ch. psittaci. Classification Because of Chlamydias unique developmental cycle, it was taxonomically classified in a separate order. Chlamydia is part of the order Chlamydiales, family Chlamydiaceae. In the early 1990s six species of Chlamydia were known. A major re-description of the Chlamydiales order in 1999, using the then new techniques of DNA analysis, split three of the species from the genus Chlamydia and reclassified them in the then newly created genus Chlamydophila, and also added three new species to this genus. In 2001 many bacteriologists strongly objected to the reclassification, although in 2006 some scientists still supported the distinctness of Chlamydophila. In 2009 the validity of Chlamydophila was challenged by newer DNA analysis techniques, leading to a proposal to "reunite the Chlamydiaceae into a single genus, Chlamydia". This appears to have been accepted by the community, bringing the number of (valid) Chlamydia species up to 9. Many probable species were subsequently isolated, but no one bothered to name them. In 2013 a 10th species was added, Ch. ibidis, known only from feral sacred ibis in France. Two more species were added in 2014 (but validated 2015): Ch. avium which infects pigeons and parrots, and Ch. gallinacea infecting chickens, guinea fowl and turkeys. Ch. abortus was added in 2015, and the Chlamydophila species reclassified. A number of new species were originally classified as aberrant strains of Ch. psittaci Genomes Chlamydia species have genomes around 1.0 to 1.3 megabases in length. Most encode ~900 to 1050 proteins. Some species also contain a DNA plasmids or phage genomes (see Table). The elementary body contains an RNA polymerase responsible for the transcription of the DNA genome after entry into the host cell cytoplasm and the initiation of the growth cycle. Ribosomes and ribosomal subunits are found in these bodies. Table 1. Genome features of selected Chlamydia species and strains. MoPn is a mouse pathogen while strain "D" is a human pathogen. About 80% of the genes in Ch. trachomatis and Ch. pneumoniae are orthologs. Adapted after Read et al. 2000 Developmental cycle Chlamydia may be found in the form of an elementary body and a reticulate body. The elementary body is the nonreplicating infectious particle that is released when infected cells rupture. It is responsible for the bacteria's ability to spread from person to person and is analogous to a spore. The elementary body may be 0.25 to 0.30 μm in diameter. This form is covered by a rigid cell wall (hence the combining form chlamyd- in the genus name). The elementary body induces its own endocytosis upon exposure to target cells. One phagolysosome usually produces an estimated 100–1000 elementary bodies. Chlamydia may also take the form of a reticulate body, which is in fact an intracytoplasmic form, highly involved in the process of replication and growth of these bacteria. The reticulate body is slightly larger than the elementary body and may reach up to 0.6 μm in diameter with a minimum of 0.5 μm. It does not have a cell wall. When stained with iodine, reticulate bodies appear as inclusions in the cell. The DNA genome, proteins, and ribosomes are retained in the reticulate body. This occurs as a result of the development cycle of the bacteria. The reticular body is basically the structure in which the chlamydial genome is transcribed into RNA, proteins are synthesized, and the DNA is replicated. The reticulate body divides by binary fission to form particles which, after synthesis of the outer cell wall, develop into new infectious elementary body progeny. The fusion lasts about three hours and the incubation period may be up to 21 days. After division, the reticulate body transforms back to the elementary form and is released by the cell by exocytosis. Studies on the growth cycle of Ch. trachomatis and Ch. psittaci in cell cultures in vitro reveal that the infectious elementary body (EB) develops into a noninfectious reticulate body (RB) within a cytoplasmic vacuole in the infected cell. After the elementary body enters the infected cell, an eclipse phase of 20 hours occurs while the infectious particle develops into a reticulate body. The yield of chlamydial elementary bodies is maximal 36 to 50 hours after infection. A histone like protein HctA and HctB play role in controlling the differentiation between the two cell types. The expression of HctA is tightly regulated and repressed by small non-coding RNA, IhtA until the late RB to EB re-differentiation. The IhtA RNA is conserved across Chlamydia species. Pathology Most chlamydial infections do not cause symptoms. Symptomatic infections often include a burning sensation when urinating and abdominal or genital pain and discomfort. All people who have engaged in sexual activity with potentially infected individuals may be offered one of several tests to diagnose the condition. Nucleic acid amplification tests (NAAT), which include polymerase chain reaction (PCR), transcription-mediated amplification (TMA), ligase chain reaction (LCR), and strand displacement amplification (SDA), are the most widely used diagnostic test for Chlamydia. Evolution Recent phylogenetic studies have revealed that Chlamydia likely shares a common ancestor with cyanobacteria, the group containing the endosymbiont ancestor to the chloroplasts of modern plants, hence, Chlamydia retains unusual plant-like traits, both genetically and physiologically. In particular, the enzyme L,L-diaminopimelate aminotransferase, which is related to lysine production in plants, is also linked with the construction of chlamydial peptidoglycan, which is required for division. The genetic encoding for the enzymes is remarkably similar in plants, cyanobacteria, and Chlamydia, demonstrating a close common ancestry. Phylogeny	Biology and health sciences	Gram-negative bacteria	Plants
2142673	https://en.wikipedia.org/wiki/Indigofera	Indigofera	Indigofera is a large genus of over 750 species of flowering plants belonging to the pea family Fabaceae. They are widely distributed throughout the tropical and subtropical regions of the world. Description Indigofera is a varied genus that has shown unique characteristics making it an interesting candidate as a potential perennial crop. Specifically, there is diverse variation among species with a number of unique characteristics. Some examples of this diversity include differences in pericarp thickness, fruit type, and flowering morphology. The unique characteristics it has displayed include potential for mixed smallholder systems with at least one other species and a resilience that allows for constant nitrogen uptake despite varying conditions. Tree Species of Indigofera are mostly shrubs, though some are small trees or herbaceous perennials or annuals. The branches are covered with silky hairs. Most of them have pinnate leaves made of three foliolates with short petioles. Small flowers grow in the leaf axils from long peduncles or spikes, their petals come in hues of red or purple, but there are a few greenish-white and yellow-flowered species. Indigofera flowers have open carpels, their organ primordial is often formed at deeper layers than other eudicots. This variety could have significant implications on its role in an actual perennial polyculture. For example, different flowering morphologies could be artificially selected for in varying directions in order to better fit in different environmental conditions and with different populations of other plants. Fruit The fruit is a long, cylindrical legume pod of varying size and shape. The types of fruit produced by different species of Indigofera can also be divided into broad categories that again show great variation. The three basic types of fruit categories can be separated by their curvature including straight, slightly curved, and falcate (sickle-shaped). In addition, several of the species including Indigofera microcarpa, Indigofera suffruticosa, and Indigofera enneaphylla have shown delayed dehiscence (maturing) of fruits This variation could again allow for artificial selection of the most abundant and nutritious fruit types and shapes. Another way to categorize Indigofera is by its pericarp thickness. The pericarp (the tissue from the ovary that surrounds the seeds) can be categorized as type I, type II, and type III with type I having the thinnest pericarp and fewest layers of schlerenchymatous (stiff) tissue and type III having the thickest pericarp and most schlerenchymatous layers. Despite the previous examples of delayed dehiscence, most fruits of this genus show normal explosive dehiscence to disperse seeds. Similar to fruit shape, the variation in fruit sizes allows for the thickest and most bountiful fruits to be selected. Uses Indigo dye Several species, especially Indigofera tinctoria and Indigofera suffruticosa, are used to produce the dye indigo. Scraps of Indigo-dyed fabric likely dyed with plants from the genus Indigofera discovered at Huaca Prieta predate Egyptian indigo-dyed fabrics by more than 1,500 years. Colonial planters in the Caribbean grew indigo and transplanted its cultivation when they settled in the colony of South Carolina and North Carolina where people of the Tuscarora confederacy adopted the dyeing process for head wraps and clothing. Exports of the crop did not expand until the mid-to late 18th century. When Eliza Lucas Pinckney and enslaved Africans successfully cultivated new strains near Charleston it became the second most important cash crop in the colony (after rice) before the American Revolution. It comprised more than one-third of all exports in value. The chemical aniline, from which many important dyes are derived, was first synthesized from Indigofera suffruticosa (syn. Indigofera anil, whence the name aniline). In Indonesia, the Sundanese use Indigofera tinctoria (known locally as tarum or nila) as dye for batik. Marco Polo was the first to report on the preparation of indigo in India. Indigo was quite often used in European easel painting during the Middle Ages. Species Indigofera comprises the following species: Palaeotropical clade Indigofera argentea Burm.f. Indigofera atriceps Hook.f. subsp. atriceps Hook.f. subsp. glandulosissima (R.E.Fr.) J.B.Gillett subsp. kaessneri (Baker f.) J.B.Gillett subsp. ramosa (Cronquist) J.B.Gillett subsp. rhodesiaca J.B.Gillett subsp. setosissima (Harms) J.B.Gillett subsp. ufipaensis J.B.Gillett Indigofera bainesii Baker Indigofera basiflora J.B.Gillett Indigofera biglandulosa J.B.Gillett Indigofera bracteolata DC. Indigofera brevicalyx Baker f. Indigofera brevipatentes Indigofera colutea (Burm.f.) Merr.—rusty indigo, sticky indigo Indigofera compressa Lam. Indigofera congesta Baker Indigofera demissa Taub. Indigofera eremophila Thulin Indigofera erythrogramma Baker Indigofera gairdnerae Baker f. Indigofera glabra L. Indigofera grata E.Mey. Indigofera hermannioides J.B.Gillett Indigofera heterotricha DC. Indigofera heudelotii Baker Indigofera hilaris Eckl. & Zeyh. var. hilaris Eckl. & Zeyh. var. microscypha (Baker) J.B. Gillett Indigofera inhambanensis Klotzsch Indigofera kirkii Oliv. Indigofera leucotricha E.Pritzel Indigofera macrocalyx Guill. & Perr. Indigofera microcalyx Baker Indigofera mildbraediana J.B.Gillett Indigofera mimosoides Baker Indigofera monantha Baker f. Indigofera mooneyi Thulin Indigofera montoya Spanish Indigo Indigofera mysorensis DC. Indigofera nebrowniana J.B.Gillett Indigofera nigritana Hook.f. Indigofera nyassica Gilli Indigofera omissa J.B.Gillett Indigofera paniculata Pers. subsp. gazensis (Baker f.) J.B.Gillett subsp. paniculata Pers. Indigofera phymatodea Thulin Indigofera podophylla Harv. Indigofera poliotes Eckl. & Zeyh. Indigofera pulchra Willd. Indigofera quarrei Cronquist Indigofera rothii Baker Indigofera rubroglandulosa Germish. Indigofera simplicifolia Lam. Indigofera strobilifera (Hochst.) Baker subsp. lanuginosa (Baker f.) J.B.Gillett subsp. strobilifera (Hochst.) Baker Indigofera suaveolens Jaub. & Spach Indigofera tanganyikensis Baker f. Indigofera tetrasperma Pers. Indigofera trachyphylla Oliv. Indigofera uniflora Roxb. Indigofera vohemarensis Baill. Indigofera wightii Wight & Arn. Indigofera wituensis Baker f. Pantropical clade Indigofera amblyantha Craib Indigofera amorphoides Jaub. & Spach Indigofera arrecta A.Rich.—Natal indigo, Bengal indigo, Java indigo Indigofera articulata Gouan Indigofera astragalina DC. Indigofera atropurpurea Hornem. Indigofera australis Willd.—Australian indigo Indigofera baumiana Harms Indigofera binderi Kotschy Indigofera blanchetiana Benth. Indigofera bojeri Baker Indigofera boranica Thulin Indigofera bosseri Du Puy & Labat Indigofera boviperda Morrison Indigofera byobiensis Hosok. Indigofera caloneura Kurz Indigofera caroliniana Mill.—Carolina indigo Indigofera cassioides DC. Indigofera cavallii Chiov. Indigofera coerulea Roxb. var. coerulea Roxb. var. monosperma (Santapau) Santapau var. occidentalis J.B.Gillett & Ali Indigofera conzattii Rose Indigofera cuernavacana Rose Indigofera cylindracea Baker Indigofera decora Lindl.—Chinese indigo var. chalara (Craib) Y.Y.Fang & C.Z.Zheng var. cooperi (Craib) Y.Y.Fang & C.Z.Zheng var. decora Lindl. var. ichangensis (Craib) Y.Y.Fang & C.Z.Zheng Indigofera deightonii J.B.Gillett subsp. deightonii J.B.Gillett subsp. rhodesica J.B.Gillett Indigofera dendroides Jacq. Indigofera dosua D.Don var. dosua D.Don var. simlensis (Ali) Sanjappa Indigofera emarginella A.Rich. Indigofera frondosa N.E.Br. Indigofera frutescens L.f. Indigofera fulgens Baker subsp. brachybotrys (Baker) J.B.Gillett subsp. fulgens Baker Indigofera galegoides DC. Indigofera georgei E.Pritz. Indigofera grandiflora B.H.Choi & S.K.Cho Indigofera haplophylla F.Muell. Indigofera hebepetala Baker var. glabra Ali var. hebepetala Baker Indigofera hedyantha Eckl. & Zeyh. Indigofera heterantha Brandis—Himalayan indigo Indigofera himalayensis Ali Indigofera hirsuta L.—hairy indigo, rough hairy indigo Indigofera homblei Baker f. & Martin Indigofera ixocarpa Peter G.Wilson & Rowe Indigofera jucunda Schrire Indigofera karnatakana Sanjappa Indigofera kirilowii Maxim ex Palib.—Kirilow's indigo Indigofera koreana Ohwi—Korean indigo Indigofera lacei Craib Indigofera langebergensis L.Bolus Indigofera laxiracemosa Baker f. Indigofera leprieurii Baker f. Indigofera longiracemosa Baill. Indigofera lyallii Baker subsp. lyallii Baker subsp. nyassica J.B.Gillett Indigofera longibarbata Engl. Indigofera longimucronata Baker f. Indigofera macrophylla Schum. & Thonn. Indigofera mangokyensis "R.Vig., p.p.A" Indigofera melanadenia Harv. Indigofera natalensis Bolus Indigofera nigrescens King & Prain Indigofera pendula Franch. var. pendula Franch. var. umbrosa (Craib) Y.Y.Fang & C.Z.Zheng Indigofera platycarpa Rose Indigofera podocarpa Baker f. & Martin Indigofera pratensis F.Muell. Indigofera rhynchocarpa Baker var. latipinna (Johnson) J.B.Gillett var. rhynchocarpa Baker var. uluguruensis J.B.Gillett Indigofera roseocaerulea Baker f. Indigofera rugosa Benth. Indigofera sanguinea N.E.Br. Indigofera schlechteri Baker f. Indigofera sedgewickiana Vatke Indigofera setiflora Baker Indigofera sokotrana Vierh. Indigofera sootepensis Craib Indigofera stenophylla Guill. & Perr. Indigofera subcorymbosa Baker Indigofera suffruticosa Mill.—anil indigo, anil de pasto Indigofera sutherlandoides Baker Indigofera swaziensis Bolus subsp. perplexa (N.E.Br.) J.B.Gillett subsp. swaziensis Bolus Indigofera thibaudiana DC. Indigofera tinctoria L.—indigo, true indigo, dye indigo subsp. arcuata (J.B.Gillett) Schrire subsp. tinctoria L. Indigofera tristis E.Mey. Indigofera truxillensis Kunth Indigofera varia E.Mey. Indigofera venulosa Benth. Indigofera verrucosa Eckl. & Zeyh. Indigofera verruculosa Peter G.Wilson Indigofera vicioides Jaub. & Spach var. rogersii R.E.Fr. var. vicioides Jaub. & Spach Indigofera zeyheri Eckl. & Zeyh. Indigofera zollingeriana Miq.—Zollinger's indigo Cape clade Indigofera alopecuroides (Burm.f.) DC. Indigofera alpina Eckl. & Zeyh. Indigofera amoena Aiton Indigofera angustata E.Mey. Indigofera angustifolia L. var. angustifolia L. var. tenuifolia (Lam.) Harv. Indigofera brachystachya (DC.) E.Mey. Indigofera burchellii DC. Indigofera candolleana Meissner Indigofera capillaris Thunb. Indigofera concava Harv. Indigofera cuneifolia Eckl. & Zeyh. Indigofera cytisoides Thunb. Indigofera declinata E.Mey. Indigofera denudata Thunb. Indigofera digitata Thunb. Indigofera dimidiata Walp. Indigofera filicaulis Eckl. & Zeyh. Indigofera filifolia Thunb. Indigofera flabellata Harv. Indigofera gifbergensis C.H.Stirt. & Jarvie Indigofera glomerata E.Mey. Indigofera heterophylla Thunb. Indigofera hispida Eckl. & Zeyh. Indigofera ionii Jarvie & C.H.Stirt. Indigofera mauritanica (L.) Thunb. Indigofera merxmuelleri A.Schreib. Indigofera meyeriana Eckl. & Zeyh. Indigofera mollis Eckl. & Zeyh. Indigofera nigromontana Eckl. & Zeyh. Indigofera nudicaulis E.Mey. Indigofera ovata Thunb. Indigofera porrecta Eckl. & Zeyh. Indigofera psoraloides (L.) L. Indigofera sarmentosa L.f. Indigofera sulcata DC. Indigofera superba C.H.Stirt. Tethyan clade Indigofera achyranthoides Taub. Indigofera alternans DC. Indigofera ammoxylum (DC.) Polhill Indigofera anabibensis A.Schreib. Indigofera angulosa Edgew. Indigofera antunesiana Harms Indigofera arabica Jaub. & Spach Indigofera aspera DC. Indigofera asperifolia Benth. var. asperifolia Benth. var. lanceolata Chodat & Hassler var. macrophylla Chodat & Hassler Indigofera auricoma E.Mey. Indigofera bemarahaensis Du Puy & Labat Indigofera bongardiana (Kuntze) Burkart Indigofera bongensis Kotschy & Peyr. Indigofera cerighellii M.Pelt. Indigofera cloiselii Drake Indigofera conjugata Baker var. conjugata Baker var. schweinfurthii (Taub.) J.B. Gillett var. trimorphophylla (Taub.) J.B. Gillett Indigofera cordifolia Roth Indigofera daleoides Harv. Indigofera dalzellii T.Cooke Indigofera depauperata Drake Indigofera depressa Harv. Indigofera dionaeifolia (S. Moore) Schrire Indigofera diphylla Vent. Indigofera disticha Eckl. & Zeyh. Indigofera diversifolia DC. Indigofera drepanocarpa Taub. Indigofera ewartiana Domin Indigofera exellii Torre Indigofera fanshawei J.B.Gillett Indigofera glandulosa Wendl. var. glandulosa Wendl. var. sykesii Baker Indigofera glaucescens Eckl. & Zeyh. Indigofera guaranitica Hassl. Indigofera gypsacea Thulin Indigofera hartwegii Rydb. Indigofera hiranensis Thulin Indigofera hochstetteri Baker subsp. hochstetteri Baker subsp. streyana (Merxm.) A.Schreib. Indigofera hololeuca Harv. Indigofera humbertiana M.Pelt. Indigofera interrupta (Du Puy, Labat & Schrire) Schrire Indigofera jamaicensis Spreng. Indigofera kelleri Baker f. Indigofera leptocarpa Eckl. & Zeyh. Indigofera leptosepala Nutt. Indigofera lespedezioides Kunth var. acutifolia Hassler var. lespedezioides Kunth Indigofera leucoclada Baker Indigofera linifolia (L.f.) Retz. Indigofera linnaei Ali Indigofera longidentata (Du Puy, Labat & Schrire) Schrire Indigofera lupatana Baker f. Indigofera mahafalensis (Du Puy, Labat & Schrire) Schrire Indigofera marmorata Balf.f. Indigofera microcarpa Desv. Indigofera miniata Ortega—coastal indigo, scarlet-pea Indigofera nephrocarpa Balf.f. Indigofera nephrocarpoides J.B.Gillett Indigofera nummularia Baker Indigofera obcordata Eckl. & Zeyh. Indigofera oblongifolia Forssk. Indigofera praticola Baker f. Indigofera pseudocompressa (Du Puy, Labat & Schrire) Schrire Indigofera pungens E.Mey. Indigofera schimperi Jaub. & Spach Indigofera semitrijuga Forssk. Indigofera sessiliflora DC. Indigofera spicata Forssk. Indigofera spiniflora Boiss. Indigofera spinosa Forssk. Indigofera squalida Prain Indigofera subulata Vahl ex Poir. Indigofera tephrosioides Kunth Indigofera thomsonii Baker f. Indigofera torulosa E.Mey. var. angustiloba (Baker f.) J.B.Gillett var. torulosa E.Mey. Indigofera trifoliata L.—threeleaf indigo var. duthiei (Naik) Sanjappa var. trifoliata L. var. unifoliolata (Merr.) De Kort & G.Thijsse Indigofera trigonelloides Jaub. & Spach Indigofera trita L.f.—Asian indigo var. maffei (Chiov.) Ali var. marginulata (Wight & Arn.) Sanjappa var. scabra (Roth) De Kort & G.Thijsse var. trita L.f. Indigofera volkensii Taub. Unassigned Indigofera acanthinocarpa Blatt. Indigofera acanthoclada Dinter Indigofera accepta N.E.Br. Indigofera acutiflora N.E.Br. Indigofera acutipetala Y.Y.Fang & C.Z.Zheng Indigofera adenocarpa E.Mey. Indigofera adenoides Baker f. Indigofera adesmiifolia A.Gray Indigofera ambelacensis Schweinf. Indigofera amitina N.E.Br. Indigofera ammobia Maconochie Indigofera ancistrocarpa Thulin Indigofera andrewsiana J.B.Gillett Indigofera andringitrensis R.Vig. Indigofera ankaratrensis R.Vig. Indigofera aquae-nitentis Bremek. Indigofera aralensis Gagnep. Indigofera arenophila Schinz Indigofera argutidens Craib Indigofera aristata Spreng. Indigofera arnottii (Kuntze) Peter G. Wilson Indigofera aspalathoides DC. Indigofera asterocalycina Gilli Indigofera atrata N.E.Br. Indigofera atricephala J.B.Gillett Indigofera baileyi F.Muell. Indigofera balfouriana Craib Indigofera bancroftii Peter G.Wilson Indigofera bangweolensis R.E.Fr. Indigofera banii N.D.Khoi & Yakovlev Indigofera barteri Hutch. & Dalziel Indigofera basedowii E.Pritz. subsp. basedowii E.Pritz. subsp. longibractea (J.Black) Peter G.Wilson Indigofera bayensis Thulin Indigofera bella Prain Indigofera benguellensis Baker Indigofera berhautiana J.B.Gillett Indigofera bijuga Walp. Indigofera blaiseae Du Puy & Labat Indigofera bogdanii J.B.Gillett Indigofera boinensis R.Vig. Indigofera brachynema J.B.Gillett Indigofera bracteata Baker Indigofera brassii Baker Indigofera brevidens Benth. var. brevidens Benth. var. uncinata Benth. Indigofera brevifilamenta J.B.Gillett Indigofera breviracemosa Torre Indigofera breviviscosa J.B.Gillett Indigofera brunoniana Wall. Indigofera buchananii Burtt Davy Indigofera bungeana Walp. Indigofera burttii Baker f. Indigofera bussei J.B.Gillett Indigofera calcicola Craib Indigofera campestris Benth. var. angustifolia M.Micheli var. campestris Benth. var. intermedia Hassler Indigofera candicans Aiton Indigofera capitata Kotschy Indigofera carlesi Craib Indigofera caudata Dunn Indigofera cecilii N.E.Br. Indigofera cedrorum Dunn Indigofera chaetodonta Franch. Indigofera charlierana Schinz var. charlierana Schinz var. lata J.B.Gillett var. scaberrima (Schinz) J.B.Gillett Indigofera chenii S.S.Chien Indigofera chevalieri Tisser. Indigofera chirensis J.B.Gillett Indigofera chuniana F.P.Metcalf Indigofera ciferrii Chiov. Indigofera cinerascens Franch. Indigofera circinella Baker f. Indigofera circinnata Harv. Indigofera cliffordiana J.B.Gillett Indigofera commixta N.E.Br. Indigofera comosa N.E.Br. Indigofera complanata Spreng. Indigofera complicata Eckl. & Zeyh. Indigofera concinna Baker Indigofera conferta J.B.Gillett Indigofera confusa Prain & Baker f. Indigofera congolensis De Wild. & T.Durand var. bongensis (Baker f.) J.B. Gillett var. congolensis De Wild. & T.Durand Indigofera constricta (Thwaites) Trimen var. constricta (Thwaites) Trimen var. deorum McVaugh Indigofera corallinosperma Torre Indigofera coronillifolia Benth. Indigofera costaricensis Benth. Indigofera crebra N.E.Br. Indigofera crotalarioides (Klotzsch) Baker Indigofera cryptantha Harv. subsp. cryptantha Harv. subsp. desmodioides (Baker) Du Puy & Labat Indigofera cuitoensis Baker f. Indigofera cuneata Oliv. Indigofera cunenensis Torre Indigofera curvata J.B.Gillett Indigofera curvirostrata Thulin Indigofera cylindrica sensu auct. Indigofera damarana Merxm. & A.Schreib. Indigofera daochengensis Y.Y.Fang & C.Z.Zheng Indigofera dasyantha Baker f. Indigofera dasycephala Baker f. Indigofera dauensis J.B.Gillett Indigofera deccanensis Sanjappa Indigofera deflersii Baker f. Indigofera dekindtii Tisser. Indigofera delagoaensis J.B.Gillett Indigofera delavayi Franch. Indigofera dembianensis (Chiov.) J.B.Gillett Indigofera densa N.E.Br. Indigofera densiflora M.Martens & Galeotti Indigofera densifructa Y.Y.Fang & C.Z.Zheng Indigofera desertorum Torre Indigofera dichroa Craib Indigofera dillwynioides Harv. Indigofera discolor Rydb. Indigofera dissitiflora Oliv. Indigofera dolichochaete Craib Indigofera dolichothyrsa Baker f. Indigofera dregeana E.Mey. Indigofera dumetorum Craib Indigofera dyeri Britten Indigofera efoliata F.Muell. Indigofera egens N.E.Br. Indigofera elandsbergensis Phillipson Indigofera elliotii (Baker f.) J.B.Gillett Indigofera elwakensis J.B.Gillett Indigofera emarginata Y.Y.Fang & C.Z.Zheng Indigofera emarginelloides J.B.Gillett Indigofera emmae De Kort & G.Thijsse Indigofera enormis N.E.Br. Indigofera erecta Thunb. Indigofera eriocarpa E.Mey. Indigofera esquirolii H.Lev. Indigofera evansiana Burtt Davy Indigofera evansii Schltr. Indigofera exigua Eckl. & Zeyh. Indigofera exilis Grierson & D.G.Long Indigofera eylesiana J.B.Gillett Indigofera faulknerae J.B.Gillett Indigofera filiformis L.f. Indigofera filipes Harv. Indigofera flavicans Baker Indigofera floribunda N.E.Br. Indigofera foliosa E.Mey. Indigofera forrestii Craib Indigofera fortunei Craib Indigofera fruticosa Rose Indigofera fulcrata Harv. Indigofera fulvopilosa Brenan Indigofera fuscosetosa Baker Indigofera galpinii N.E.Br. Indigofera gangetica Sanjappa Indigofera garckeana Vatke Indigofera geminata Baker Indigofera giessii A.Schreib. Indigofera glaucifolia Cronquist Indigofera gloriosa Cronquist Indigofera goetzei Harms Indigofera gracilis Spreng. Indigofera graniticola J.B.Gillett Indigofera griseoides Harms Indigofera grisophylla Fourc. Indigofera guatemalensis Moc., Sessé & Cerv. ex Backer Indigofera guthriei Bolus Indigofera hamiltonii Duthie & Prain Indigofera hamulosa Schltr. Indigofera hancockii Craib Indigofera hantamensis Diels Indigofera helmsii Peter G.Wilson Indigofera hendecaphylla Jacq.—trailing indigo, creeping indigo, spicate indigo var. hendecaphylla Jacq. var. siamensis (Hosseus) Gagnep. Indigofera henryi Craib Indigofera heterocarpa Baker Indigofera hewittii Baker f. Indigofera hinanensis H.T.Tsai & T.F.Yu Indigofera hofmanniana Schinz Indigofera holstii (Baker f.) Baker f. Indigofera holubii N.E.Br. Indigofera howellii Craib & W.W.Sm. Indigofera huillensis Baker f. Indigofera humifusa Eckl. & Zeyh. Indigofera humilis Kunth Indigofera hundtii Rossberg Indigofera hybrida N.E.Br. Indigofera hygrobia Malme Indigofera imerinensis Du Puy & Labat Indigofera incana Thunb. Indigofera incompta McVaugh Indigofera ingrata N.E.Br. Indigofera insularis Chiov. Indigofera intermedia Harv. Indigofera intricata Boiss. Indigofera inyangana N.E.Br. Indigofera irodoensis Du Puy & Labat Indigofera ischnoclada Harms Indigofera itremoensis Du Puy & Labat Indigofera jaliscensis Rose Indigofera jikongensis Y.Y.Fang & C.Z.Zheng Indigofera jindongensis Y.Y.Fang & C.Z.Zheng Indigofera karkarensis (Thulin) Schrire Indigofera kasinii Boonyam. Indigofera kerrii De Kort & G.Thijsse Indigofera kerstingii Harms Indigofera knoblecheri Kotschy Indigofera kongwaensis J.B.Gillett Indigofera krookii Zahlbr. Indigofera kuntzei Harms Indigofera kurtzii Kuntze Indigofera lamellata Thulin Indigofera lancifolia Rydb. Indigofera lasiantha Desv. Indigofera latifolia Micheli Indigofera latisepala J.B.Gillett Indigofera laxiflora Craib Indigofera leendertziae N.E.Br. Indigofera lenticellata Craib Indigofera lepida N.E.Br. Indigofera leptoclada Harms Indigofera letestui Tisser. Indigofera leucotricha E.Pritz. Indigofera limosa L.Bolus Indigofera lindheimeriana Scheele—Lindheimer's indigo Indigofera litoralis Chun & T.Chen Indigofera livingstoniana J.B.Gillett Indigofera longicauda Thuan Indigofera longipedicellata J.B.Gillett Indigofera longipedunculata Y.Y.Fang & C.Z.Zheng Indigofera longistaminata Schrire Indigofera lotononoides Baker f. Indigofera lughensis Thulin Indigofera luzonensis De Kort & G.Thijsse Indigofera lydenburgensis N.E.Br. Indigofera macrantha Harms Indigofera madagascariensis Vatke Indigofera malacostachys Harv. Indigofera malindiensis J.B.Gillett Indigofera malongensis Cronquist Indigofera manyoniensis Baker f. Indigofera maritima Baker Indigofera masaiensis J.B.Gillett Indigofera masonae N.E.Br. Indigofera matudae Lundell Indigofera maymyoensis Sanjappa Indigofera megacephala J.B.Gillett Indigofera mekongensis Jessup Indigofera mendesii Torre Indigofera mendoncae J.B.Gillett Indigofera mengtzeana Craib Indigofera micheliana Rose Indigofera micrantha E.Mey. Indigofera micropetala Baker f. Indigofera mildrediana Torre Indigofera milne-redheadii J.B.Gillett Indigofera minbuensis Gage Indigofera mischocarpa Schltr. Indigofera mollicoma N.E.Br. Indigofera monanthoides J.B.Gillett Indigofera monbeigii Craib Indigofera monophylla DC. Indigofera monostachya Eckl. & Zeyh. Indigofera montana Rose Indigofera mouroundavensis Baill. Indigofera muliensis Y.Y.Fang & C.Z.Zheng Indigofera mundtiana Eckl. & Zeyh. Indigofera mupensis Torre subsp. abercornensis J.B.Gillett subsp. mupensis Torre Indigofera mwanzae J.B.Gillett Indigofera myosurus Craib Indigofera nairobiensis Baker f. subsp. nairobiensis Baker f. subsp. viscida J.B.Gillett Indigofera nambalensis Harms Indigofera neoglabra Wang & T.Tang Indigofera neosericopetala P.C. Li Indigofera nesophila Lievens & Urbatsch Indigofera nigricans Pers. Indigofera nivea R.Vig. Indigofera nugalensis Thulin Indigofera nummulariifolia (L.) Alston Indigofera obscura N.E.Br. Indigofera ogadensis J.B.Gillett Indigofera oligophylla Klotzsch Indigofera omariana J.B.Gillett Indigofera ormocarpoides Baker Indigofera orthocarpa C.Presl Indigofera oubanguiensis Tisser. Indigofera ovina Harv. Indigofera oxalidea Baker Indigofera oxytropis Harv. Indigofera oxytropoides Schltr. Indigofera palmeri S.Watson Indigofera pampaniniana Craib Indigofera panamensis Rydb. Indigofera pappei Fourc. Indigofera paracapitata J.B.Gillett Indigofera paraglaucifolia Torre Indigofera paraoxalidea Torre Indigofera parkesii Craib Indigofera parodiana Burkart Indigofera parviflora F. Heyne ex Hook. & Arn. Indigofera patula Baker Indigofera pauciflora Eckl. & Zeyh. Indigofera paucifolioides Blatt. & Hallb. Indigofera paucistrigosa J.B.Gillett Indigofera pearsonii Baker f. Indigofera pechuelii Kuntze Indigofera pedicellata Wight & Arn. Indigofera pedunculata Baker Indigofera pellucida J.B.Gillett & Thulin Indigofera peltata J.B.Gillett Indigofera peltieri Du Puy & Labat Indigofera penduloides Y.Y.Fang & C.Z.Zheng Indigofera perriniana Spreng. Indigofera petiolata Cronquist Indigofera phyllanthoides Baker Indigofera pilgeriana Schltr. Indigofera pilosa Poir.—softhairy indigo Indigofera pinifolia Baker Indigofera placida N.E.Br. Indigofera platypoda E.Mey. Indigofera pobeguinii J.B.Gillett Indigofera polygaloides M.B.Scott Indigofera polysphaera Baker Indigofera pongolana N.E.Br. Indigofera porrigens Colla Indigofera prieureana Guill. & Perr. Indigofera procumbens L. Indigofera prostrata Willd. Indigofera pruinosa Baker Indigofera pseudoevansii Hilliard & B.L.Burtt Indigofera pseudointricata J.B.Gillett Indigofera pseudoparvula R.Vig. Indigofera pseudoreticulata Grierson & D.G.Long Indigofera pseudosubulata Baker f. Indigofera pseudotinctoria Matsum. Indigofera pueblensis Rydb. Indigofera purpusii Brandegee Indigofera quinquefolia E.Mey. Indigofera radicifera Cronquist Indigofera ramosissima J.B.Gillett Indigofera ramulosissima Hosok. Indigofera rautanenii Baker f. Indigofera reducta N.E.Br. Indigofera rehmannii Baker f. Indigofera remota Baker f. Indigofera repens Cronquist Indigofera reticulata Franch. Indigofera retusa N.E.Br. Indigofera rhodantha Fourc. Indigofera rhytidocarpa Harv. subsp. angolensis J.B.Gillett subsp. rhytidocarpa Harv. Indigofera rigioclada Craib Indigofera ripae N.E.Br. Indigofera rojasii Hassl. Indigofera rostrata Bolus Indigofera ruspolii Baker f. Indigofera sabulosa Thulin Indigofera salmoniflora Rose Indigofera salteri Baker f. Indigofera santapaui Sanjappa Indigofera santosii Torre Indigofera saxicola Benth. Indigofera scabrida Dunn Indigofera scarciesii Scott-Elliot Indigofera schinzii N.E.Br. Indigofera schliebenii Harms Indigofera schultziana F.Muell. Indigofera scopiformis Thulin Indigofera sebungweensis J.B.Gillett Indigofera secundiflora Poir. Indigofera senegalensis Lam. Indigofera sensitiva Franch. Indigofera sericovexilla C.T.White Indigofera sesquipedalis Sanjappa Indigofera sessilifolia DC. Indigofera setosa N.E.Br. Indigofera sieberiana Scheele Indigofera silvestrii Pamp. var. alii Sanjappa var. silvestrii Pamp. Indigofera simaoensis Y.Y.Fang & C.Z.Zheng Indigofera sisalis J.B.Gillett Indigofera smutsii J.B.Gillett Indigofera sordida Harv. Indigofera souliei Craib Indigofera sparsa Baker Indigofera sparteola Chiov. Indigofera sphaerocarpa A.Gray—Sonoran indigo Indigofera sphinctosperma Standl. Indigofera splendens Ficalho & Hiern Indigofera stachyodes Lindl. Indigofera stenosepala Baker Indigofera sticta Craib Indigofera stricta L.f. Indigofera strigulosa Baker f. Indigofera suarezensis Du Puy & Labat Indigofera subargentea De Wild. Indigofera subsecunda Gagnep. Indigofera subulifera Baker Indigofera subverticellata Gagnep. Indigofera szechuensis Craib Indigofera taborensis J.B.Gillett Indigofera tanaensis J.B.Gillett Indigofera taruffiana Torre Indigofera taylori J.B.Gillett Indigofera teixeirae Torre Indigofera tengyuehensis H.T.Tsai & T.F.Yu Indigofera tenuifolia Lam. Indigofera tenuipes Polhill Indigofera tenuis Milne-Redh. subsp. major J.B.Gillett subsp. tenuis Milne-Redh. Indigofera tenuissima E.Mey. Indigofera terminalis Baker Indigofera tetraptera Taub. Indigofera texana Buckley Indigofera thesioides Jarvie & C.H.Stirt. Indigofera thikaensis J.B.Gillett Indigofera thothathri Sanjappa Indigofera thymoides Baker Indigofera tirunelvelica Sanjappa Indigofera tomentosa Eckl. & Zeyh. Indigofera torrei J.B.Gillett Indigofera transvaalensis Baker f. Indigofera trialata A.Chev. Indigofera trichopoda Guill. & Perr. Indigofera trifolioides Baker f. Indigofera triquetra E.Mey. Indigofera tristoides N.E.Br. Indigofera tryonii Domin Indigofera tumidula Rose Indigofera ufipaensis J.B.Gillett Indigofera ugandensis Baker f. Indigofera vanderystii J.B.Gillett Indigofera velutina E.Mey. Indigofera venusta Eckl. & Zeyh. Indigofera viscidissima Baker subsp. orientalis J.B.Gillett subsp. viscidissima Baker Indigofera vivax Schrank Indigofera wildiana J.B.Gillett Indigofera williamsonii (Harv.) N.E.Br. Indigofera wilsonii Craib Indigofera woodii Bolus Indigofera zanzibarica J.B.Gillett Indigofera zavattarii Chiov. Indigofera zenkeri Baker f. Indigofera zornioides Du Puy & Labat Species names with uncertain taxonomic status The status of the following species is unresolved: Indigofera abyssinica Hochst. ex Baker Indigofera adaochengensis Y.Y. Fang & C.Z. Zheng Indigofera adenophylla Graham Indigofera adenotricha Peter G.Wilson Indigofera adesmiaefolia A. Gray Indigofera adonensis E.Mey. Indigofera aeruginis Schweinf. Indigofera agowensis Hochst. ex Baker Indigofera alata Schweinf. Indigofera alba Gouault Indigofera amaliae Domin Indigofera angulata Lindl. Indigofera angulata Rottler ex Spreng. Indigofera aphylla Breiter Indigofera arborescens Zuccagni Indigofera arghawan Royle Indigofera argyrea Chiov. Indigofera armata Wall. Indigofera ascendens Walp. Indigofera astragaloides Welw. ex Romariz Indigofera athrophylla Eckl. & Zeyh. Indigofera axillaris E.Mey. Indigofera bagshawei Baker f. Indigofera baoulensis A.Chev. Indigofera barbata Desv. Indigofera barcensis Chiov. Indigofera bequaerti De Wild. Indigofera berteroana Spreng. Indigofera bertolonii Steud. Indigofera biflora Roth Indigofera bilabiata Loisel. ex Drapiez Indigofera boylei Hort. ex Vilmorin's Indigofera brachycarpa Graham Indigofera brachyodon Domin Indigofera brachyphylla Al-Turki Indigofera brachypoda Steud. ex A.Rich. Indigofera brevipes (S. Watson) Rydb. Indigofera bufalina Lour. Indigofera caesia Zipp. ex Span. Indigofera caespitosa Wight Indigofera calva E.Mey. Indigofera carlesii Craib Indigofera ceciliae N.E.Br. Indigofera celebica Miq. Indigofera centrota Eckl. & Zeyh. Indigofera chitralensis Sanjappa Indigofera cinericolor F.Muell. Indigofera clitorioides G.Don Indigofera colorata Roxb. ex Wight & Arn. Indigofera coluteifolia Jaub. & Spach Indigofera condensata De Wild. Indigofera conradsii Baker f. Indigofera constricta Rydb. Indigofera cornezuelo Moc. & Sessé ex DC. Indigofera cornuligera Peter G.Wilson & Rowe Indigofera coronillaefolia A. Cunn. ex Benth. Indigofera coronillaefolia hort. Indigofera crassisiliqua Steud. Indigofera dalzelliana (Kuntze) Peter G.Wilson Indigofera dalzielii Hutch. Indigofera debilis Graham Indigofera decumbens Hill Indigofera deginensis Sanjappa Indigofera dequinensis Sanjappa Indigofera dewevrei Micheli Indigofera diffusa Desv. Indigofera dimorphophylla Schinz Indigofera disjuncta J. B. Gillett Indigofera dodecaphylla Ficalho & Hiern Indigofera dorycnium Fenzl Indigofera dosycnium Fenzl Indigofera dubia Steud. Indigofera dumosa E.Mey. Indigofera elachantha Peter G.Wilson & Rowe Indigofera elatior Carrière Indigofera elegans Schumach. & Thonn. Indigofera ellenbeckii Baker f. Indigofera elskensii Baker f. Indigofera enonensis E.Mey. Indigofera erectifructa Y.Endo, H.Ohashi & Madulid Indigofera erythrantha Hochst. ex Baker Indigofera erythrogrammoides De Wild. Indigofera esquirolii H. Lév. Indigofera faberi Craib Indigofera flavovirens R.E.Fr. Indigofera flexuosa Eckl. & Zeyh. Indigofera flexuosa Graham Indigofera florida E.Mey. Indigofera foliolosa Graham Indigofera formosana Matsum. Indigofera franchetii X.F.Gao & Schrire Indigofera frumentacea Roxb. ex Wight & Arn. Indigofera fruticulosa Walp. Indigofera fuzi Sieb. ex Miq. Indigofera gilletii De Wild. & T.Durand Indigofera glauca Lam. Indigofera glauca Perr. ex DC. Indigofera grahamiana Steud. Indigofera grandifoliola Carrière Indigofera graveolens Schrad. Indigofera griquana Schltr. ex Zahlbr. Indigofera guineensis Schumach. & Thonn. Indigofera haematica Peter G.Wilson Indigofera hainanensis H.T.Tsai & T.T.Yü Indigofera heptaphylla Hiern Indigofera hislopii Baker f. Indigofera hockii De Wild. & Baker f. Indigofera hookeriana Meisn. Indigofera hover Forssk. Indigofera inconspicua Domin Indigofera iwafusi Sieb. ex Lavallee Indigofera jaubertiana Schweinf. Indigofera jirahulia Buch.-Ham. Indigofera juncea Decne. Indigofera karongensis Baker Indigofera kerensis Chiov. Indigofera kisantuensis De Wild. & T.Durand Indigofera kotoensis Hayata Indigofera latibracteata Harms Indigofera latipinna I.M.Johnst. Indigofera laxeracemosa Baker f. Indigofera leptocaulis Eckl. & Zeyh. Indigofera leptophylla E.Mey. Indigofera lignosa De Wild. Indigofera limifolia Benth. Indigofera lindleyana Spreng. ex Steud. Indigofera linearis DC. Indigofera linearis Guill. & Perr. Indigofera litoralis Chun & T.C. Chen Indigofera liukiuennsis Makino & Nemoto Indigofera lonchocarpifolia Baker Indigofera longebarbata Engl. Indigofera longepedicellata J. B. Gillett Indigofera longeracemosa Boivin ex Baill. Indigofera longibractea J.M.Black Indigofera lupulina Baker Indigofera machaerocarpa Fenzl ex Baker Indigofera macroptera hort. ex Lavallée Indigofera macrostachys Vent. Indigofera mangokyensis R. Vig. Indigofera marginata Walp. Indigofera masukuensis Baker Indigofera mckinlayi F.Muell. Indigofera mearnsi Standl. Indigofera megaphylla X.F.Gao Indigofera melanotricha Steud. ex A.Rich. Indigofera melolobioides Benth. ex Harv. Indigofera microphylla Lam. Indigofera microstachya C.Presl Indigofera minutiflora Hochst. ex Chiov. Indigofera minutiflora Walp. Indigofera moeroensis De Wild. Indigofera multijuga Baker Indigofera mutisii (Kunth) Spreng. Indigofera nematopoda Baker f. Indigofera neoarborea Hu ex F.T. Wang & Tang Indigofera noldeae Rossbach Indigofera nuda G.Don Indigofera nyikensis Baker Indigofera oligantha Harms ex Baker f. Indigofera oligosperma DC. Indigofera orixensis Roxb. ex Wight & Arn. Indigofera oroboides E.Mey. Indigofera oxyrachis Peter G.Wilson Indigofera paludosa Lepr. ex Guill. & Perr. Indigofera palustris Vatke Indigofera perrottetii DC. Indigofera petraea Peter G.Wilson & Rowe Indigofera pilifera Peter G.Wilson & Rowe Indigofera platyspira J.B.Gillett ex Thulin & M.G.Gilbert Indigofera plumosa Spreng. Indigofera polyclada Peter G.Wilson & Rowe Indigofera polysperma De Wild. & T.Durand Indigofera pratensis var. coriacea Domin Indigofera preladoi Harms Indigofera pretoriana Harms Indigofera procumbens Torre Indigofera propinqua Hochst. ex Chiov. Indigofera psammophila Peter G.Wilson Indigofera pseudoheterantha X.F.Gao & Schrire Indigofera pseudomoniliformis Schrire Indigofera purpurea Page ex Steud. Indigofera quadrangularis Graham Indigofera racemosa L. Indigofera rarifolia Steud. Indigofera rechodes Eckl. & Zeyh. Indigofera reflexa E.Mey. Indigofera rhechodes Walp. Indigofera rhodosantha Zipp. ex Miq. Indigofera rigescens E.Mey. Indigofera roylei Koehne Indigofera roylii Hort. ex Dippel Indigofera rubromarginata Thulin Indigofera rumphiensis Schrire Indigofera rupestris Eckl. & Zeyh. Indigofera rupicola Peter G.Wilson & Rowe Indigofera sabulicola Benth. Indigofera saltiana Steud. Indigofera sangana Harms, in Schltr. Indigofera scabrella Kazandj. & Peter G.Wilson Indigofera schimperiana Hochst. Indigofera scoparia Vahl ex DC. Indigofera secunda E.Mey. Indigofera sericea Benth. ex Baker Indigofera sericea L. Indigofera sericea Thunb. ex Harv. Indigofera sericophylla Franch. Indigofera setacea E.Mey. Indigofera shipingensis X.F.Gao Indigofera shirensis Taub. ex Baker f. Indigofera signata Domin Indigofera similis N.E.Br. Indigofera sinuspersica Mozaff. Indigofera socotrana Vierh. Indigofera sofa Scott-Elliot Indigofera solirimae Schrire Indigofera somalensis Vatke Indigofera sousae M.A.Exell Indigofera sparsiflora Hochst. ex Baker Indigofera speciosa Fraser ex Hook. Indigofera spirocarpa Harms Indigofera spoliata Hoffmanns. Indigofera subincana N.E.Br. Indigofera subquadriflora Hochst. ex Chiov. Indigofera subtilis E.Mey. Indigofera sylvatica Sieber ex Spreng. Indigofera sylvestris Pamp. Indigofera taiwaniana T.C.Huang & M.J.Wu Indigofera tenella Schumach. & Thonn. Indigofera tenella Vahl ex DC. Indigofera tenuicaulis Klotzsch Indigofera tenuisiliqua Schweinf. Indigofera ternata Roxb. ex Wight & Arn. Indigofera thirionni H.Lév. Indigofera thonningii Schumach. & Thonn. Indigofera tinctaria Hook. Indigofera triflora Peter G.Wilson & Rowe Indigofera trita var. nubica (J.B.Gillett) L.Boulos & Schrire Indigofera tritoidea Baker Indigofera ultima (Kuntze) Peter G.Wilson Indigofera unifoliata Merr. Indigofera urostachya Fenzl ex Baker Indigofera viguieri Callm. & Labat Indigofera villosa Berg. ex Walp. Indigofera wannanii Peter G.Wilson Indigofera wentzeliana Harms Indigofera wynbergensis S.Moore Indigofera zig-zag De Wild. Ecology Indigofera species are used as food plants by the larvae of some Lepidoptera species, including the turnip moth (Agrotis segetum).	Biology and health sciences	Fabales	Plants
2143933	https://en.wikipedia.org/wiki/Rhinorrhea	Rhinorrhea	Rhinorrhea (American English), also spelled rhinorrhoea or rhinorrhœa (British English), or informally runny nose is the free discharge of a thin mucus fluid from the nose; it is a common condition. It is a common symptom of allergies (hay fever) or certain viral infections, such as the common cold or COVID-19. It can be a side effect of crying, exposure to cold temperatures, cocaine abuse, or drug withdrawal, such as from methadone or other opioids. Treatment for rhinorrhea may be aimed at reducing symptoms or treating underlying causes. Rhinorrhea usually resolves without intervention, but may require treatment by a doctor if symptoms last more than 10 days or if symptoms are the result of foreign bodies in the nose. The term rhinorrhea was coined in 1866 from the Greek rhino- ("of the nose") and -rhoia ("discharge" or "flow"). Signs and symptoms Rhinorrhea is characterized by an excess amount of mucus produced by the mucous membranes that line the nasal cavities. The membranes create mucus faster than it can be processed, causing a backup of mucus in the nasal cavities. As the cavity fills up, it blocks off the air passageway, causing difficulty breathing through the nose. Air caught in nasal cavities – namely the sinus cavities, cannot be released and the resulting pressure may cause a headache or facial pain. If the sinus passage remains blocked, there is a chance that sinusitis may result. If the mucus backs up through the Eustachian tube, it may result in ear pain or an ear infection. Excess mucus accumulating in the throat or back of the nose may cause a post-nasal drip, resulting in a sore throat or coughing. Additional symptoms include sneezing, nosebleeds, and nasal discharge. Causes A runny nose can be caused by anything that irritates or inflames the nasal tissues, including infections such as the common cold and influenza, and allergies and various irritants. Some people have a chronically runny nose for no apparent reason (non-allergic rhinitis or vasomotor rhinitis). Less common causes include polyps, a foreign body, a tumor or migraine-like headaches. Some causes of rhinorrhea include: acute sinusitis (nasal and sinus infection), allergies, chronic sinusitis, common cold, coronaviruses (COVID-19), decongestant nasal spray overuse, deviated septum, dry air, eosinophilic granulomatosis with polyangiitis, granulomatosis with polyangiitis, hormonal changes, influenza (flu), lodged object, medicines (such as those used to treat high blood pressure, erectile dysfunction, depression, seizures and other conditions), nasal polyps, non-allergic rhinitis (chronic congestion or sneezing not related to allergies), occupational asthma, pregnancy, respiratory syncytial virus (RSV), spinal fluid leak, and tobacco smoke. Cold temperatures Rhinorrhea is especially common in cold weather. Cold-induced rhinorrhea occurs due to a combination of thermodynamics and the body's natural reactions to cold weather stimuli. One of the purposes of nasal mucus is to warm inhaled air to body temperature as it enters the body; this requires the nasal cavities to be constantly coated with liquid mucus. In cold weather the mucus lining nasal passages tends to dry out, so that mucous membranes must work harder, producing more mucus to keep the cavity lined. As a result, the nasal cavity can fill up with mucus. At the same time, when air is exhaled, water vapor in breath condenses as the warm air meets the colder outside temperature near the nostrils. This causes excess water to build up inside nasal cavities, spilling out through the nostrils. Inflammatory Infection Rhinorrhea can be a symptom of other diseases, such as the common cold or influenza. During these infections, the nasal mucous membranes produce excess mucus, filling the nasal cavities. This is to prevent infection from spreading to the lungs and respiratory tract, where it could cause far worse damage. It has also been suggested that viral rhinorrhea is a result of viral evolution whereby virus variants that increase nasal secretion and are thus more resistant to the body's immune defenses are selected for. Rhinorrhea caused by these infections usually occur on circadian rhythms. Over the course of a viral infection, sinusitis (the inflammation of the nasal tissue) may occur, causing the mucous membranes to release more mucus. Acute sinusitis consists of the nasal passages swelling during a viral infection. Chronic sinusitis occurs when sinusitis continues for longer than three months. Allergies Rhinorrhea can also occur when individuals with allergies to certain substances, such as pollen, dust, latex, soy, shellfish, or animal dander, are exposed to these allergens. In people with sensitized immune systems, the inhalation of one of these substances triggers the production of the antibody immunoglobulin E (IgE), which binds to mast cells and basophils. IgE bound to mast cells are stimulated by pollen and dust, causing the release of inflammatory mediators such as histamine. In the nasal cavities, these inflammatory mediators cause inflammation and swelling of the tissue, as well as increased mucus production. Particulate matter in polluted air and chemicals such as chlorine and detergents, which can normally be tolerated, can make the condition considerably worse. Crying Rhinorrhea is also associated with shedding tears (lacrimation), whether from emotional events or from eye irritation. When excess tears are produced, the liquid drains through the inner corner of the eyelids, through the nasolacrimal duct, and into the nasal cavities. As more tears are shed, more liquid flows into the nasal cavities, both stimulating mucus production and hydrating any dry mucus already present in the nasal cavity. The buildup of fluid is usually resolved via mucus expulsion through the nostrils. Non-inflammatory Head trauma Rhinorrhea can be caused by a head injury, a serious condition. A basilar skull fracture can result in a rupture of the barrier between the sinonasal cavity and the anterior cranial fossae or the middle cranial fossae. This can cause the nasal cavity to fill with cerebrospinal fluid (cerebrospinal fluid rhinorrhoea, CSF rhinorrhea), a condition that can lead to a number of serious complications, including death if not addressed properly. Other causes Rhinorrhea can occur as a symptom of opioid withdrawal accompanied by lacrimation. Other causes include cystic fibrosis, whooping cough, nasal tumors, hormonal changes, and cluster headaches. Rhinorrhea can also be the side effect of several genetic disorders, such as primary ciliary dyskinesia, as well as common irritants such as spicy foods, nail polish remover, or paint fumes. Treatment In most cases, treatment for rhinorrhea is not necessary since it will clear up on its own, especially if it is the symptom of an infection. For general cases nose-blowing can get rid of the mucus buildup. Though blowing may be a quick-fix solution, it increases mucosal production in the sinuses, leading to frequent and higher mucus buildups in the nose in the medium term. Alternatively, saline or vasoconstrictor nasal sprays may be used, but may become counterproductive after several days of use, causing rhinitis medicamentosa. In some cases, such as those due to allergies or sinus infections, there are medicinal treatments available. Several types of antihistamines can be obtained relatively cheaply to treat cases caused by allergies; antibiotics may help in cases of bacterial sinus infections.	Biology and health sciences	Symptoms and signs	Health
2144540	https://en.wikipedia.org/wiki/Toxaphene	Toxaphene	Toxaphene was an insecticide used primarily for cotton in the southern United States during the late 1960s and the 1970s. Toxaphene is a mixture of over 670 different chemicals and is produced by reacting chlorine gas with camphene. It can be most commonly found as a yellow to amber waxy solid. Toxaphene was banned in the United States in 1990 and was banned globally by the 2001 Stockholm Convention on Persistent Organic Pollutants. It is a very persistent chemical that can remain in the environment for 1–14 years without degrading, particularly in the soil. Testing performed on animals, mostly rats and mice, has demonstrated that toxaphene is harmful to animals. Exposure to toxaphene has proven to stimulate the central nervous system, as well as induce morphological changes in the thyroid, liver, and kidneys. Toxaphene has been shown to cause adverse health effects in humans. The main sources of exposure are through food, drinking water, breathing contaminated air, and direct contact with contaminated soil. Exposure to high levels of toxaphene can cause damage to the lungs, nervous system, liver, kidneys, and in extreme cases, may even cause death. It is thought to be a potential carcinogen in humans, though this has not yet been proven. Composition Toxaphene is a synthetic organic mixture composed of over 670 chemicals, formed by the chlorination of camphene (C10H16) to an overall chlorine content of 67–69% by weight. The bulk of the compounds (mostly chlorobornanes, chlorocamphenes, and other bicyclic chloroorganic compounds) found in toxaphene have chemical formulas ranging from C10H11Cl5 to C10H6Cl12, with a mean formula of C10H10Cl8. The formula weights of these compounds range from 308 to 551 grams/mole; the theoretical mean formula has a value of 414 grams/mole. Toxaphene is usually seen as a yellow to amber waxy solid with a piney odor. It is highly insoluble in water but freely soluble in aromatic hydrocarbons and readily soluble in aliphatic organic solvents. It is stable at room temperature and pressure. It is volatile enough to be transported for long distances through the atmosphere. Applications Advertisements for Toxaphene were seen in agricultural periodicals such as Farm Journal as early as 1950. Toxaphene was primarily used as a pesticide for cotton in the southern United States during the late 1960s and 1970s. It was also used on small grains, maize, vegetables, and soybeans. Outside of the realm of crops, it was also used to control ectoparasites such as lice, flies, ticks, mange, and scam mites on livestock. In some cases it was used to kill undesirable fish species in lakes and streams. The breakdown of usage can be summarized: 85% on cotton, 7% to control insect pests on livestock and poultry, 5% on other field crops, 3% on soybeans, and less than 1% on sorghum. The first recorded usage of toxaphene was in 1966 in the United States, and by the early to mid 1970s, toxaphene was the United States' most heavily used pesticide. Over 34 million pounds of toxaphene were used annually from 1966 to 1976. As a result of Environmental Protection Agency restrictions, annual toxaphene usage fell to 6.6 million pounds in 1982. In 1990, the EPA banned all usage of toxaphene in the United States. Toxaphene is still used in countries outside the United States but much of this usage has been undocumented. Between 1970 and 1995, global usage of toxaphene was estimated to be 670 million kilograms (1.5 billion pounds). Production Toxaphene was first produced in the United States in 1947 although it was not heavily used until 1966. By 1975, toxaphene production reached its peak at 59.4 million pounds annually. Production decreased more than 90% from this value by 1982 due to Environmental Protection Agency restrictions. Overall, an estimated 234,000 metric tons (over 500 million pounds) have been produced in the United States. Between 25% and 35% of the toxaphene produced in the United States has been exported. There are currently 11 toxaphene suppliers worldwide. Environmental effects When released into the environment, toxaphene can be quite persistent and exists in the air, soil, and water. In water, it can evaporate easily and is fairly insoluble. Its solubility is 3 mg/L of water at 22 degrees Celsius. Toxaphene breaks down very slowly and has a half-life of up to 12 years in the soil. It is most commonly found in air, soil, and sediment found at the bottom of lakes or streams. It can also be present in many parts of the world where it was never used because toxaphene is able to evaporate and travel long distances through air currents. Toxaphene can eventually be degraded, through dechlorination, in the air using sunlight to break it down. The degradation of toxaphene usually occurs under aerobic conditions. The levels of toxaphene have decreased since its ban. However, due to its persistence, it can still be found in the environment today. Exposure The three main paths of exposure to toxaphene are ingestion, inhalation, and absorption. For humans, the main source of toxaphene exposure is through ingested seafood. When toxaphene enters the body, it usually accumulates in fatty tissues. It is broken down through dechlorination and oxidation in the liver, and the byproducts are eliminated through feces. People that live near an area that has high toxaphene contamination are at high risk to toxaphene exposure through inhalation of contaminated air or direct skin contact with contaminated soil or water. Eating large quantities of fish on a daily basis also increases susceptibility to toxaphene exposure. Finally, exposure is rare, yet possible through drinking water when contaminated by toxaphene runoff from the soil. However, toxaphene has been rarely seen at high levels in drinking water due to toxaphene's nearly complete insolubility in water. Shellfish, algae, fish and marine mammals have all been shown to exhibit high levels of toxaphene. People in the Canadian Arctic, where a traditional diet consists of fish and marine animals, have been shown to consume ten times the accepted daily intake of toxaphene. Also, blubber from beluga whales in the Arctic were found to have unhealthy and toxic levels of toxaphene. Health effects In humans When inhaled or ingested, sufficient quantities of toxaphene can damage the lungs, nervous system, and kidneys, and may cause death. The major health effects of toxaphene involve central nervous system stimulation leading to convulsive seizures. The dose necessary to induce nonfatal convulsions in humans is about 10 milligrams per kilogram body weight per day. Several deaths linked to toxaphene have been documented in which an unknown quantity of toxaphene was ingested intentionally or accidentally from food contamination. The deaths are attributed to respiratory failure resulting from seizures. Chronic inhalation exposure in humans results in reversible respiratory toxicity. A study conducted between 1954 and 1972 of male agricultural workers and agronomists exposed to toxaphene and other pesticides showed that there are higher proportions of bronchial carcinoma in the test group than in the unexposed general population. However, toxaphene may not have been the main pesticide responsible for tumor production. Tests on lab animals show that toxaphene causes liver and kidney cancer, so the EPA has classified it as a Group B2 carcinogen, meaning it is a probable human carcinogen. The International Agency for Research on Cancer has classified it as a Group 2B carcinogen. Toxaphene can be detected in blood, urine, breast milk, and body tissues if a person has been exposed to high levels, but it is removed from the body quickly, so detection has to occur within several days of exposure. It is not known whether toxaphene can affect reproduction in humans. In animals Toxaphene was used to treat mange in cattle in California in the 1970s and there were reports of cattle deaths following the toxaphene treatment. Chronic oral exposure in animals affects the liver, the kidney, the spleen, the adrenal and thyroid glands, the central nervous system, and the immune system. Toxaphene stimulates the central nervous system by antagonizing neurons leading to hyperpolarization of neurons and increased neuronal activity. Regulations Toxaphene has been found on at least 68 of the 1,699 National Priorities List sites identified by the United States Environmental Protection Agency. Toxaphene has been forbidden in Germany since 1980. Most uses of toxaphene were cancelled in the U.S. in 1982 with the exception of use on livestock in emergency situations, and for controlling insects on banana and pineapple crops in Puerto Rico and the U.S. Virgin Islands. All uses of toxaphene were cancelled in the U.S. in 1990. Toxaphene has been banned in 37 countries, including Austria, Belize, Brazil, Costa Rica, Dominican Republic, Egypt, the EU, India, Ireland, Kenya, Korea, Mexico, Panama, Singapore, Thailand and Tonga. Its use has been severely restricted in 11 other countries, including Argentina, Columbia, Dominica, Honduras, Nicaragua, Pakistan, South Africa, Turkey, and Venezuela. In the Stockholm Convention on POPs, which came into effect on 17 May 2004, twelve POPs were listed to be eliminated or their production and use restricted. The OCPs or pesticide-POPs identified on this list have been termed the "dirty dozen" and include aldrin, chlordane, DDT, dieldrin, endrin, heptachlor, hexachlorobenzene, mirex, and toxaphene. The EPA has determined that lifetime exposure to 0.01 milligrams per liter of toxaphene in the drinking water is not expected to cause any adverse noncancer effects if the only source of exposure is drinking water, and has established the maximum contaminant level (MCL) of toxaphene at 0.003 mg/L. The United States Food and Drug Administration (FDA) uses the same level for the maximum level permissible in bottled water. The FDA has determined that the concentration of toxaphene in bottled drinking water should not exceed 0.003 milligrams per liter. The United States Department of Transportation lists toxaphene as a hazardous material and has special requirements for marking, labeling, and transporting the material. It is classified as an extremely hazardous substance in the United States as defined in Section 302 of the U.S. Emergency Planning and Community Right-to-Know Act (42 U.S.C. 11002), and is subject to strict reporting requirements by facilities which produce, store, or use it in significant quantities. Trade names Trade names and synonyms include Chlorinated camphene, Octachlorocamphene, Camphochlor, Agricide Maggot Killer, Alltex, Allotox, Crestoxo, Compound 3956, Estonox, Fasco-Terpene, Geniphene, Hercules 3956, M5055, Melipax, Motox, Penphene, Phenacide, Phenatox, Strobane-T, Toxadust, Toxakil, Vertac 90%, Toxon 63, Attac, Anatox, Royal Brand Bean Tox 82, Cotton Tox MP82, Security Tox-Sol-6, Security Tox-MP cotton spray, Security Motox 63 cotton spray, Agro-Chem Brand Torbidan 28, and Dr Roger's TOXENE.	Technology	Pest and disease control	null
15865266	https://en.wikipedia.org/wiki/Waru%20Waru	Waru Waru	Waru Waru is an Aymara term for the agricultural technique developed by pre-Hispanic people in the Andes region of South America from Ecuador to Bolivia; this regional agricultural technique is also referred to as camellones in Spanish. Functionally similar agricultural techniques have been developed in other parts of the world, all of which fall under the broad category of raised field agriculture. This type of altiplano field agriculture consists of parallel canals alternated by raised planting beds, which would be strategically located on floodplains or near a water source so that the fields could be properly irrigated. These flooded fields were composed of soil that was rich in nutrients due to the presence of aquatic plants and other organic materials. Through the process of mounding up this soil to create planting beds, natural, recyclable fertilizer was made available in a region where nitrogen-rich soils were rare. By trapping solar radiation during the day, this raised field agricultural method also protected crops from freezing overnight. These raised planting beds were irrigated very efficiently by the adjacent canals which extended the growing season significantly, allowing for more food yield. Waru Waru were able to yield larger amounts of food than previous agricultural methods due to the overall efficiency of the system. This technique is dated to around 300 B.C., and is most commonly associated with the Tiwanaku culture of the Lake Titicaca region in southern Bolivia, who used this method to grow crops like potatoes and quinoa. This type of agriculture also created artificial ecosystems, which attracted other food sources such as fish and lake birds. Past cultures in the Lake Titicaca region likely utilized these additional resources as a subsistence method. It combines raised beds with irrigation channels to prevent damage by soil erosion during floods. These fields ensure both collecting of water (either fluvial water, rainwater or phreatic water) and subsequent drainage. The drainage aspect of this method makes it particularly useful in many areas subjected to risks of brutal floods, such as tropical parts of Bolivia and Peru where it emerged. Raised field agricultural methods have been used in many other countries such as China, Mexico and Belize. Mexican Chinampas were similar to Waru Waru in that they were created on or near a water source in order to properly irrigate crops. Raised fields are known in Belize from various sites, including Pulltrouser Swamp. Modern Uses In the 1960s, geographers William Denevan, George Plafker, and Kenneth Lee found evidence of raised-field agriculture that had been utilized in the Llanos de Moxos region of Bolivia's Amazon basin, a region that was previously thought to have been unable to sustain large-scale agriculture because of what was believed to have been an unfavorable rainforest environment. This discovery led to a joint experimental archaeology project in the region involving archaeologist Clark Erickson, the Inter-American Foundation, the Parroquia of San Ignacio, the Bolivian Institute of Archaeology, and the University of Pennsylvania Museum of Archaeology and Anthropology. The goal of this experiment was to attempt to restore indigenous raised-field agriculture in the region. This project began in 1990 at the Biological Station of the Beni Department in Bolivia. Because of the experiment's success, it was later implemented further in collaboration with local indigenous communities. The indigenous community provided land for the project and the Inter-American Foundation paid them wages to build and maintain the plots, which successfully produced manioc and maize. These plots did not require extensive upkeep following the initial season's planting, and were self-sufficient because of the artificial ecosystems that they created. This agricultural method was also revived by Alan Kolata of the University of Chicago in 1984, in Tiwanaku, Bolivia as well as Puno, Peru. Research on Waru Waru and its effectiveness in the past has led to a resurgence of the technique amongst contemporary Aymara- and Quechua-speaking native peoples in Bolivia and Peru. By utilizing this centuries-old technique, modern people in the region have been able to make use of the harsh altiplano landscape around Lake Titicaca. This method is now being used in different areas of South America where farming is difficult, such as the altiplano and the Amazon basin. Because of this method, indigenous people are now able to farm the landscape much more efficiently and without the use of modern equipment. This method also allows for large-scale agriculture to be performed in the Amazon basin without having to rely on deforestation. Experiments Research was done at two raised-field sites by Diego Sanchez de Lozada et al. in the northern altiplano of Bolivia near Lake Titicaca in an effort to better understand the effects of frost on potato crops. At an altitude of , these crops were subject to temperature and moisture variation. Temperatures of the soil on top of the high raised mounds was about 1 degree Celsius higher than the temperature of the ground in nearby fields, showing that the raised-field technique was able to partially mitigate frost effects on potato crops at night. Temperature and moisture analysis of the raised fields showed that the higher temperature present was due to above-ground processes, which caused cold air to fall to the canals and not on the planted rows. The frost mitigation effects of the raised field system kept crops from freezing overnight, which increased crop yield. History Lake Titicaca Region 16th Century Spanish accounts of the Lake Titicaca region mentioned the different types of agriculture utilized by the native peoples in detail, however there was never any mention of raised fields in their records. The lack of Spanish accounts strongly suggests that these Waru Waru were no longer in use by the time the conquistadors reached the Lake Titicaca region. The raised fields of the region are numerous and range in size, however they are generally wide, long, and tall. These pre-Hispanic fields cover about of land in Bolivia and Peru, and sit above an altitude of around 3,800 m. Radiocarbon dates taken from habitation sites associated with raised field agriculture in the region indicate usage sometimes between 1000 B.C. to A.D. 400. Thermoluminescence dating was also used to date pottery shards in associated areas, the results of which agree with the radiocarbon dates. Field stratigraphy was used to provide relative dates of the usage of certain raised fields in the area. The habitation sites in association with these fields indicate large populations and long-term occupations, suggesting that raised field agriculture was able to sustain large numbers of people. These dates provided from Andean sites suggest that this form of agriculture was a relatively early phenomenon in the area that slowly expanded throughout the region, and was utilized by various cultures during different time periods.	Technology	Buildings and infrastructure	null
15868921	https://en.wikipedia.org/wiki/Amp%C3%A8re%27s%20force%20law	Ampère's force law	In magnetostatics, the force of attraction or repulsion between two current-carrying wires (see first figure below) is often called Ampère's force law. The physical origin of this force is that each wire generates a magnetic field, following the Biot–Savart law, and the other wire experiences a magnetic force as a consequence, following the Lorentz force law. Equation Special case: Two straight parallel wires The best-known and simplest example of Ampère's force law, which underlaid (before 20 May 2019) the definition of the ampere, the SI unit of electric current, states that the magnetic force per unit length between two straight parallel conductors is where is the magnetic force constant from the Biot–Savart law, is the total force on either wire per unit length of the shorter (the longer is approximated as infinitely long relative to the shorter), is the distance between the two wires, and , are the direct currents carried by the wires. This is a good approximation if one wire is sufficiently longer than the other, so that it can be approximated as infinitely long, and if the distance between the wires is small compared to their lengths (so that the one infinite-wire approximation holds), but large compared to their diameters (so that they may also be approximated as infinitely thin lines). The value of depends upon the system of units chosen, and the value of decides how large the unit of current will be. In the SI system, with the magnetic constant, in SI units General case The general formulation of the magnetic force for arbitrary geometries is based on iterated line integrals and combines the Biot–Savart law and Lorentz force in one equation as shown below. where is the total magnetic force felt by wire 1 due to wire 2 (usually measured in newtons), and are the currents running through wires 1 and 2, respectively (usually measured in amperes), The double line integration sums the force upon each element of wire 1 due to the magnetic field of each element of wire 2, and are infinitesimal vectors associated with wire 1 and wire 2 respectively (usually measured in metres); see line integral for a detailed definition, The vector is the unit vector pointing from the differential element on wire 2 towards the differential element on wire 1, and \|r\| is the distance separating these elements, The multiplication × is a vector cross product, The sign of is relative to the orientation (for example, if points in the direction of conventional current, then ). To determine the force between wires in a material medium, the magnetic constant is replaced by the actual permeability of the medium. For the case of two separate closed wires, the law can be rewritten in the following equivalent way by expanding the vector triple product and applying Stokes' theorem: In this form, it is immediately obvious that the force on wire 1 due to wire 2 is equal and opposite the force on wire 2 due to wire 1, in accordance with Newton's third law of motion. Historical background The form of Ampere's force law commonly given was derived by James Clerk Maxwell in 1873 and is one of several expressions consistent with the original experiments of André-Marie Ampère and Carl Friedrich Gauss. The x-component of the force between two linear currents I and I, as depicted in the adjacent diagram, was given by Ampère in 1825 and Gauss in 1833 as follows: Following Ampère, a number of scientists, including Wilhelm Weber, Rudolf Clausius, Maxwell, Bernhard Riemann, Hermann Grassmann, and Walther Ritz, developed this expression to find a fundamental expression of the force. Through differentiation, it can be shown that: and also the identity: With these expressions, Ampère's force law can be expressed as: Using the identities: and Ampère's results can be expressed in the form: As Maxwell noted, terms can be added to this expression, which are derivatives of a function Q(r) and, when integrated, cancel each other out. Thus, Maxwell gave "the most general form consistent with the experimental facts" for the force on ds arising from the action of ds': Q is a function of r, according to Maxwell, which "cannot be determined, without assumptions of some kind, from experiments in which the active current forms a closed circuit." Taking the function Q(r) to be of the form: We obtain the general expression for the force exerted on ds by ds''' : Integrating around s' eliminates k and the original expression given by Ampère and Gauss is obtained. Thus, as far as the original Ampère experiments are concerned, the value of k has no significance. Ampère took k=−1; Gauss took k=+1, as did Grassmann and Clausius, although Clausius omitted the S component. In the non-ethereal electron theories, Weber took k=−1 and Riemann took k=+1. Ritz left k undetermined in his theory. If we take k = −1, we obtain the Ampère expression: If we take k=+1, we obtain Using the vector identity for the triple cross product, we may express this result as When integrated around ds' the second term is zero, and thus we find the form of Ampère's force law given by Maxwell: Derivation of parallel straight wire case from general formula Start from the general formula: Assume wire 2 is along the x-axis, and wire 1 is at y=D, z=0, parallel to the x-axis. Let be the x-coordinate of the differential element of wire 1 and wire 2, respectively. In other words, the differential element of wire 1 is at and the differential element of wire 2 is at . By properties of line integrals, and . Also, and Therefore, the integral is Evaluating the cross-product: Next, we integrate from to : If wire 1 is also infinite, the integral diverges, because the total attractive force between two infinite parallel wires is infinity. In fact, what we really want to know is the attractive force per unit length'' of wire 1. Therefore, assume wire 1 has a large but finite length . Then the force vector felt by wire 1 is: As expected, the force that the wire feels is proportional to its length. The force per unit length is: The direction of the force is along the y-axis, representing wire 1 getting pulled towards wire 2 if the currents are parallel, as expected. The magnitude of the force per unit length agrees with the expression for shown above. Notable derivations Chronologically ordered: Ampère's original 1823 derivation: Maxwell's 1873 derivation: Treatise on Electricity and Magnetism vol. 2, part 4, ch. 2 (§§502–527) Pierre Duhem's 1892 derivation: (EPUB) translation of: Leçons sur l'électricité et le magnétisme vol. 3, appendix to book 14, pp. 309-332 Alfred O'Rahilly's 1938 derivation: Electromagnetic Theory: A Critical Examination of Fundamentals vol. 1, pp. 102–104 (cf. the following pages, too)	Physical sciences	Magnetostatics	Physics
9205088	https://en.wikipedia.org/wiki/Pelagornithidae	Pelagornithidae	The Pelagornithidae, commonly called pelagornithids, pseudodontorns, bony-toothed birds, false-toothed birds or pseudotooth birds, are a prehistoric family of large seabirds. Their fossil remains have been found all over the world in rocks dating between the Early Paleocene and the Pliocene-Pleistocene boundary. Most of the common names refer to these birds' most notable trait: tooth-like points on their beak's edges, which, unlike true teeth, contained Volkmann's canals and were outgrowths of the premaxillary and mandibular bones. Even "small" species of pseudotooth birds were the size of albatrosses; the largest ones had wingspans estimated at 5–6 metres (15–20 ft) and were among the largest flying birds ever to live. They were the dominant seabirds of most oceans throughout most of the Cenozoic, and modern humans apparently missed encountering them only by a tiny measure of evolutionary time: the last known pelagornithids were contemporaries of Homo habilis and the beginning of the history of technology. Description and ecology The biggest of the pseudotooth birds were the largest flying birds known. Almost all of their remains from the Neogene are immense, but in the Paleogene there were a number of pelagornithids that were around the size of a great albatross (genus Diomedea) or even a bit smaller. The undescribed species provisionally called "Odontoptila inexpectata" – from the Paleocene-Eocene boundary of Morocco – is the smallest pseudotooth bird discovered to date and was just a bit larger than a white-chinned petrel (Procellaria aequinoctialis). The Pelagornithidae had extremely thin-walled bones widely pneumatized with the air sac extensions of the lungs. Most limb bone fossils are very much crushed for that reason. In life, the thin bones and extensive pneumatization enabled the birds to achieve large size while remaining below critical wing loadings. Though 25 kg/m2 (5 lb/ft2) is regarded as the maximum wing loading for powered bird flight, there is evidence that bony-toothed birds used dynamic soaring flight almost exclusively: the proximal end of the humerus had an elongated diagonal shape that could hardly have allowed for the movement necessary for the typical flapping flight of birds; their weight thus cannot be easily estimated. The attachment positions for the muscles responsible for holding the upper arm straightly outstretched were particularly well-developed, and altogether the anatomy seems to allow for an ability of holding the wings rigidly at the glenoid joint unmatched by any other known bird. This is especially prominent in the Neogene pelagornithids, and less developed in the older Paleogene forms. The sternum had the deep and short shape typical of dynamic soarers, and bony outgrowths at the keel's forward margin securely anchored the furcula. The legs were proportionally short, the feet probably webbed and the hallux was vestigial or entirely absent; the tarsometatarsi (anklebones) resembled those of albatrosses while the arrangement of the front toes was more like in fulmars. Typical for pseudotooth birds was a second toe that attached a bit kneewards from the others and was noticeably angled outwards. The "teeth" were probably covered by the rhamphotheca in life, and there are two furrows running along the underside of the upper bill just inside the ridges which bore the "teeth". Thus, when the bill was closed only the upper jaw's "teeth" were visible, with the lower ones hidden behind them. Inside the eye sockets of at least some pseudotooth birds – perhaps only in the younger species – were well-developed salt glands. Altogether, almost no major body part of pelagornithids is known from a well-preserved associated fossil and most well-preserved material consists of single bones only; on the other hand the long occurrence and large size makes for a few rather comprehensive (though much crushed and distorted) remains of individual birds that were entombed by as they lay dead, complete with some fossilized feathers. Large parts of the skull and some beak pieces are found not too infrequently. In February 2009, an almost-complete fossilized skull of a presumed Odontopteryx from around the Chasicoan-Huayquerian boundary c. 9 million years ago (Ma) was unveiled in Lima. It had been found a few months earlier in Ocucaje District of Ica Province, Peru. According to paleontologist Mario Urbina, who discovered the specimen, and his colleagues Rodolfo Salas, Ken Campbell and Daniel T. Ksepka, the Ocucaje skull is the best-preserved pelagornithid cranium known as of 2009. Ecology and extinction Unlike the true teeth of Mesozoic stem-birds like Archaeopteryx or Ichthyornis, the pseudoteeth of the pelagornithids do not seem to have had serrated or otherwise specialized cutting edges, and were useful to hold prey for swallowing whole rather than to tear bits off it. Since the teeth were hollow or at best full of cancellous bone and are easily worn or broken off in fossils, it is surmised they were not extremely resilient in life either. Pelagornithid prey would thus have been soft-bodied, and have encompassed mainly cephalopods and soft-skinned fishes. Prey items may have reached considerable size. Though some reconstructions show pelagornithids as diving birds in the manner of gannets, the thin-walled highly pneumatized bones which must have fractured easily judging from the state of fossil specimens make such a mode of feeding unlikely, if not outright dangerous. Rather, prey would have been picked up from immediately below the ocean surface while the birds were flying or swimming, and they probably submerged only the beak in most situations. Their quadrate bone articulation with the lower jaw resembled that of a pelican or other birds that can open their beak widely. Altogether, the pseudotooth birds would have filled an ecological niche almost identical to that of the larger fish-eating pteranodontian pterosaurs, whose extinction at the end of the Cretaceous may well have paved the way for the highly successful 50-million-year reign of the Pelagornithidae. Like them as well as modern albatrosses, the pseudotooth birds could have used the system of ocean currents and atmospheric circulation to take round-track routes soaring over the open oceans, returning to breed only every few years. Unlike albatrosses today, which avoid the tropical equatorial currents with their doldrums, Pelagornithidae were found in all sorts of climates, and records from around 40 Ma stretch from Belgium through Togo to the Antarctic. It is conspicuous that penguins and plotopterids – both wing-propelled divers that foraged over the continental shelf – are almost invariably found in the company of pseudotooth birds. Thus, pseudotooth birds seem to have gathered in some numbers in upwelling regions, presumably to feed but perhaps also to breed nearby. It is sometimes claimed that as with some other seabirds (e.g. the flightless Plotopteridae), the evolutionary radiation of cetaceans and pinnipeds outcompeted the pseudotooth birds and drove them into extinction. While this may be correct for the plotopterids, for pelagornithids it is not so likely for two reasons: First, the Pelagornithidae continued to thrive for 10 million years after modern-type baleen whales evolved, and in the Middle Miocene Pelagornis coexisted with Aglaocetus and Harrison's whale (Eobalaenoptera harrisoni) in the Atlantic off the Eastern Seaboard, while the Pacific Osteodontornis inhabited the same seas as Balaenula and Morenocetus; the ancestral smallish sperm whale genus Aulophyseter (and/or Orycterocetus) occurred in both Northern Hemisphere oceans at that time, while the mid-sized sperm whale Brygmophyseter roamed the North Pacific. As regards Miocene pinnipeds, a diversity of ancient walruses and ancestral fur seals like Thalassoleon inhabited the north-east, while the ancient leopard seal Acrophoca is a remarkable species known from the south-east Pacific. Secondly, pinnipeds are limited to near-shore waters while pseudotooth birds roamed the seas far and wide, like large cetaceans, and like all big carnivores all three groups were K-strategists with moderate to very low population densities. Thus, direct competition for food between bony-toothed birds and cetaceans or pinnipeds cannot have been very severe. As both the birds and pinnipeds would need level ground near the sea to raise their young, competition for breeding grounds may have affected the birds' population. In that respect, the specializations for dynamic soaring restricted the number of possible nesting sites for the birds, but on the other hand upland on islands or in coastal ranges could have provided breeding grounds for Pelagornithidae that was inaccessible for pinnipeds; just as many albatrosses today nest in the uplands of islands (e.g. the Galápagos or Torishima). The bony-toothed birds probably required strong updrafts for takeoff and would have preferred higher sites anyway for this reason, rendering competition with pinniped rookeries quite minimal. As regards breeding grounds, giant eggshell fragments from the Famara mountains on Lanzarote, Canary Islands, were tentatively attributed to Late Miocene pseudotooth birds. As regards the Ypresian London Clay of the Isle of Sheppey, wherein pelagornithid fossils are not infrequently found, it was deposited in a shallow epicontinental sea during a very hot time with high sea levels. The presumed breeding sites cannot have been as far offshore as many seabird rookeries are today, as the region was hemmed in between the Alps and the Grampian and Scandinavian Mountains, in a sea less wide than the Caribbean is today. Neogene pseudotooth birds are common along the America coasts near the Appalachian and Cordilleran mountains, and these species thus presumably also bred not far offshore or even in the mountains themselves. In that respect the presence of medullary bone in the specimens from Lee Creek Mine in North Carolina, United States, is notable, as among birds this is generally only found in laying females, indicating that the breeding grounds were probably not far away. At least Pacific islands of volcanic origin would be eroded away in the last millions of years however, obliterating any remains of pelagornithid breeding colonies that might have once existed in the open ocean. Necker Island for example was of significant size 10 million years ago, when Osteodontornis roamed the Pacific. There is no obvious single reason for the pseudotooth birds' extinction. A scenario of general ecological change – exacerbated by the beginning ice age and changes in ocean currents due to plate tectonic shifts (e.g. the emergence of the Antarctic circumpolar current or the closing of the Isthmus of Panama) – is more likely, with the pseudotooth birds as remnants of the world's Paleogene fauna ultimately failing to adapt. In that respect it may be significant that some lineages of cetaceans, like the primitive dolphins of the Kentriodontidae or the shark-toothed whales, flourished contemporary with the Pelagornithidae and became extinct at about the same time. Also, the modern diversity of pinniped and cetacean genera evolved largely around the Mio-Pliocene boundary, suggesting that many ecological niches emerged or became vacant. In addition, whatever caused the Middle Miocene disruption and the Messinian Salinity Crisis did affect the trophic web of Earth's oceans not insignificantly either, and the latter event led to a widespread extinction of seabirds. Together, this combination of factors led to Neogene animals finally replacing the last remnants of the Paleogene fauna in the Pliocene. In that respect, it is conspicuous that the older pseudotooth birds are typically found in the same deposits as plotopterids and penguins, while younger forms were sympatric with auks, albatrosses, penguins and Procellariidae – which, however, underwent an adaptive radiation of considerable extent coincident (and probably caused by) with the final demise of the Paleogene-type trophic web. Although the fossil record is necessarily incomplete, as it seems cormorants and seagulls were very rarely found in association with the Pelagornithidae. Irrespective of the cause of their ultimate extinction, during the long time of their existence the pseudotooth birds furnished prey for large predators themselves. Few if any birds that coexisted with them were large enough to harm them while airborne; as evidenced by the Early Eocene Limnofregata, the frigatebirds coevolved with the Pelagornithidae and may well have harassed any of the small species for food on occasion, as they today harass albatrosses. From the Middle Miocene or Early Pliocene of the Lee Creek Mine, some remains of pseudotooth birds which probably fell victim to sharks while feeding are known. The large members of the abundant Lee Creek Mine shark fauna that hunted near the water's surface included the broadnose sevengill shark (Notorynchus cepedianus), Carcharias sand tiger sharks, Isurus and Cosmopolitodus mako sharks, Carcharodon white sharks, the snaggletooth shark Hemipristis serra, tiger sharks (Galeocerdo), Carcharhinus whaler sharks, the lemon shark (Negaprion brevirostris) and hammerhead sharks (Sphyrna), and perhaps (depending on the bird fossils' age) also Pristis sawfishes, Odontaspis sand tiger sharks, and Lamna and Parotodus benedeni mackerel sharks. It is notable that fossils of smaller diving birds – for example auks, loons and cormorants – as well as those of albatrosses are much more commonly found in those shark pellets than pseudotooth birds, supporting the assumption that the latter had quite low population densities and caught much of their food in mid-flight. A study on Pelagornis flight performance suggests that, unlike modern seabirds, it relied on thermal soaring much like continental soaring birds and Pteranodon. External appearance Nothing is known for sure about the colouration of these birds, as they left no living descendants. But some inferences can be made based on their phylogeny: if they were a member of the "higher waterbird" group (see below), they most probably had a plumage similar to that depicted in the reconstruction of Osteodontornis orri – Procellariiformes and Pelecaniformes in the modern sense (or Ciconiiformes, if Pelecaniformes are merged there) have hardly any carotenoid or structural colors at all in their plumage, and generally lack even phaeomelanins. Thus, the only colours commonly found in these birds are black, white and various shades of grey. Some have patches of iridescent feathers, or brownish or reddish hues, but these are rare and limited in extent, and those species in which they are found (e.g. bitterns, ibises or the hammerkop) are generally only found in freshwater habitat. If the pseudotooth birds are Galloanseres, phaeomelanins might be more likely to have occurred in their feathers, but it is notable that the most basal lineages of Anseriformes are typically grey-and-black or black-and-white. Among ocean-going birds in general, the upperside tends to be much darker than the underside (including the underwings) – though some petrels are dark grey all over, a combination of more or less dark grey upperside and white underside and (usually) head is a widespread colouration found in seabirds and may either be plesiomorphic for "higher waterbirds" or, perhaps more likely, be an adaptation to provide camouflage, in particular against being silhouetted against the sky if seen by prey in the sea. It is notable that at least the primary remiges, and often the other flight feathers too, are typically black in birds – even if the entire remaining plumage is completely white, as in some pelicans or in the Bali starling (Leucopsar rothschildi). This is due to the fact that melanins will polymerize, making all-black feathers very robust; as the largest forces encountered by bird feathers affect the flight feathers, the large amount of melanin gives them better resistance against being damaged in flight. In soaring birds as dependent on strong winds as the bony-toothed birds were, black wingtips and perhaps tails can be expected to have been present. As regards the bare parts, all the presumed close relatives of the Pelagornithidae quite often have rather bright reddish colours, in particular on the bill. The phylogenetic uncertainties surrounding them do not allow to infer whether the bony-toothed birds had a throat sac similar to pelicans. If they did, it was probably red or orange, and may have been used in mating displays. Sexual dimorphism was probably almost nonexistent, as it typically is among the basal Anseriformes and the "higher waterbirds". Taxonomy, systematics and evolution The name "pseudodontorns" refers to the genus Pseudodontornis, which for some time served as the family's namesake. However, the presently used name Pelagornithidae pre-dates Pseudodontornithidae, and thus modern authors generally prefer "pelagornithids" over "pseudodontorns". The latter name is generally found in mid-20th-century literature however. Historically, the disparate bones of the pseudotooth birds were spread across six groups: a number of genera described from leg bones was placed in a family Cyphornithidae, and considered close allies of the pelican family (Pelecanidae). They were united with the latter in a superfamily Pelecanides in suborder Pelecanae, or later on (after the endings of taxonomic ranks were fixed to today's standard) Pelecanoidea in suborder Pelecani. Subsequently, some allied them with the entirely spurious "family" "Cladornithidae" in a "pelecaniform" suborder "Cladornithes". Those genera known from skull material were typically assigned to one or two families (Odontopterygidae and sometimes also Pseudodontornithidae) in a "pelecaniform" suborder Odontopteryges or Odontopterygia. Pelagornis meanwhile, described from wing bones, was traditionally placed in a monotypic "pelecaniform" family Pelagornithidae. This was often assigned either to the gannet and cormorant suborder Sulae (which was formerly treated as superfamily Sulides in suborder Pelecanae), or to the Odontopterygia. The sternum of Gigantornis was placed in the albatross family (Diomedeidae) in the order of tube-nosed seabirds (Procellariiformes). The most extensive taxonomic and systematic confusion affected Dasornis. That genus was established based on a huge skull piece, which for long was placed in the Gastornithidae merely due to its size. Argillornis – nowadays recognized to belong in Dasornis – was described from wing bones, and generally included in the Sulae as part of the "Elopterygidae" – yet another invalid "family", and its type genus is generally not considered a modern-type bird by current authors. Some additional tarsometatarsus (ankle) bone fragments were placed in the genus Neptuniavis and assigned to the Procellariidae in the Procellariiformes. All these remains were only shown to belong in the pseudotooth bird genus Dasornis in 2008. The most basal known pelagornithid is Protodontopteryx. Systematics and phylogeny The systematics of bony-toothed birds are subject of considerable debate. Initially, they were allied with the (then-polyphyletic) "Pelecaniformes" (pelicans and presumed allies, such as gannets and frigatebirds) and the Procellariiformes (tube-nosed seabirds like albatrosses and petrels), because of their similar general anatomy. Some of the first remains of the massive Dasornis were mistaken for a ratite and later a gastornithid. They were even used to argue for a close relationship between these two groups – and indeed, the pelicans and tubenoses, as well as for example the other "Pelecaniformes" (cormorants and allies) which are preferably separated as Phalacrocoraciformes nowadays, the Ciconiiformes (storks and/or either herons and ibises or the "core" Pelecaniformes) and Gaviiformes (loons/divers) seem to make up a radiation, possibly a clade, of "higher waterbirds". However, the Pelagornithidae are not generally held to be a missing link between pelicans and albatrosses anymore, but if anything much closer to the former and only convergent to the latter in ecomorphology. In 2005, a cladistic analysis proposed a close relationship between pseudotooth birds and waterfowl (Anseriformes). These are not part of the "higher waterbirds" but of the Galloanserae, a basal lineage of neognath birds. Some features, mainly of the skull, support this hypothesis. For example, the pelagornithids lack a crest on the underside of the palatine bone, while the Neoaves – the sister clade of the Galloanserae which includes the "higher waterbirds" and the "higher landbirds" – have such a crest. Also, like ducks, geese and swans pelagornithids only have two and not three condyles on the mandibular process of the quadrate bone, with the middle condyle beakwards of the side condyle. Their basipterygoid articulation is similar to that of the Galloanseres. At the side of the parasphenoid lamina, there is a wide platform as in Anseriformes. The bony-toothed birds' attachment of the coronoideal part of the external mandible adductor muscles was located at the midline, the rostropterygoid process had a support at its base and the mesethmoid bone had a deep depression for the caudal concha, just as in waterfowl. As regards other parts of the skeleton, the proposed synapomorphies of pelagornithids and waterfowl are found mainly in the arm- and handbones: the ulna had a strongly convex upper backside at its elbow end – at the handward end of which the scapulotricipital muscles attached –, a point-tipped dorsal cotyle and only a shallow depression to house the meniscus between ulna and radius; towards the elbow, the intercondylar sulcus of the ulna becomes wide and is bordered by a long winding ridge on the belly side. The radius, meanwhile, has a convex ventral border to the humeral cotyle, which prominently continues the hind edge of the knob where the biceps brachii muscle attaches; towards the upper side of the radius bone the surface becomes flat and triangular handwards of the articular surface for the ulna. The carpometacarpus of both Anseriformes and pseudotooth birds has a prominent pisiform process, which extends from the carpal trochlea far fingerwards along the bone's forward side. On the carpometacarpus' underside, there is a long but narrow symphysis of the distal metacarpals, with the large metacarpal bone having a mid-ridge that at its outer end curves tailwards, and the thumb joint has a well-developed knob on the hind side of its articular surface. The leg and foot bones, as is to be expected from birds not as specialized for swimming as waterfowl are, show less similarities between Anseriformes and pseudotooth birds: on the tibiotarsus there is a wide incision between the condyles and the middle condyle is narrower than the side condyle and protrudes forwards; the tarsometatarsus has a low distal vascular foramen with recessed opening on its plantar surface and a middle toe trochlea that is elongated, slightly oblique, projects to the underside of the foot and is pointed at the tip. It is unclear what to make of these apomorphies supposedly uniting Anseriformes and bony-toothed birds, for on the other hand, the sternum, distal humerus, leg and foot bones of pelagornithids seem to show apomorphies typical of "higher waterbirds". While details of the braincase bones are held to be very informative phylogenetically, the skull features in which the two groups are similar are generally related to the point where the bill attaches to the skull, and thus might have been subject to the selective forces brought about by skimming food from the upper water layer. The apparent non-neoavian traits distinguishing pelagornithids could just as well be retained or atavistic plesiomorphies; as the "higher waterbirds" are very ancient Neoaves and none of the suspected basal members of their radiation (see also "Graculavidae") were included in the analysis, it is not known for sure when the derived conditions typical of modern Neoaves were acquired. Footbone traits are notoriously prone to selection forces in birds, with convergent evolution known to inhibit or even invalidate cladistic analyses; however, the apparent autapomorphies of the lower arm and hand bones are hard to explain by anything else than an actual relationship. The location of the salt glands inside the eye sockets of Osteodontornis, Pelagornis (and probably others) shows that whatever their relationships were, the pelagornithids adapted to an oceanic habitat independently from penguins and tubenoses, which instead have supraorbital salt glands. Their missing or vestigial hallux – like in ducks but unlike in pelicans which have all four toes fully developed and webbed – was held against a close relationship with pelicans. But as is known today, pelicans are closer to storks (which have a hallux but no webbing) than to pseudotooth birds and evolved their fully webbed toes independently. With both a webbed and a hypotrophied hallux being apomorphic and paraphyletic, its absence in pseudotooth birds does not provide much information on their relationship. While giant Galloanserae were common and diverse in the Paleogene in particular, these (Gastornis and mihirungs) were flightless terrestrial birds; it is perhaps significant though that the only other "bone-toothed" birds known so far are the two species of the moa-nalo genus Thambetochen, extinct giant flightless dabbling ducks from the Hawaiian Islands. In any case, the 2005 cladistic analysis uses a representative sample of Procellariiformes and recovers them as strongly supported clade in agreement with the current consensus. The presumed close relationship between bony-toothed birds and tubenoses can thus be disregarded after all. As regards "Pelecaniformes", the analysis does not recover the correct phylogeny and does not include the shoebill (Balaeniceps rex, a "missing link" between pelicans and storks) either; clearly, the adaptive radiation of the pelican-stork lineage is misleading the analysis here. In addition, the Galloanserae are not recovered as monophyletic. In 2007, a far more comprehensive cladistic analysis of bird anatomy including some fossil forms (though not the crucial Late Cretaceous taxa, which are usually known only from fragmentary remains) resolved the "higher waterbird" radiation somewhat better; still, the problem of leg and foot traits confounding the analysis was noticeable. As their relationships are still unresolved between Galloanserae and "higher waterbirds", the pseudotooth birds are here placed in the distinct order Odontopterygiformes as a compromise, rather than in a pelecaniform/ciconiiform or anseriform suborder Odontopterygia or even a family of the Anseriformes, Ciconiiformes or Pelecaniformes. Such a treatment is unlikely to be completely wrong in either case, as the pseudotooth birds are well distinct from the Presbyornithidae and Scopidae, today generally regarded as the very basal divergences of, respectively, the Anseriformes and the pelican-stork group. It also provides leeway should the proposed separation of the Pelagornithidae into several families turn out to be appropriate. It is perhaps notable that when Boris Spulski established the Odontopterygia in 1910, he did this partly because he noted some of the similarities between pseudotooth birds and waterfowl listed above. Dasornis was long mistaken for a gastornithid, now strongly suspected to be very close indeed to the Anseriformes. Also, the pelagornithid Palaeochenoides mioceanus was initially mistaken for an anseriform, and the same might hold true for the supposed Oligocene swan Guguschia nailiae. In the former case, however, a "much the more convincing" analysis for a placement outside the Galloanseres was published the year after its description already. Most unrecognised pelagornithid bones were initially assigned to "higher waterbird" families however, typically to the (then-paraphyletic) "Pelecaniformes", but in particular the tarsometatarsus was typically mistaken for that of a procellariiform. The Odontopterygiformes were first proposed when Osteodontornis was described from the first – and still the only known – reasonably complete skeleton of one of these birds. Hildegarde Howard found that, no matter that some of its features resembled other birds, the combination was quite unlike any neognath known. While the authors claim it is beyond the paper's scope, the study describing Protodontopteryx suggests that the proposed pro-galloansere traits might actually be plesiomorphic in relation to Aves. It also notes "striking" similarities between pelagornithids and Ichthyornis in terms of jaw anatomy, but still classifies them as neognaths due to the well-developed hypotarsal crests, a supratendineal bridge on the distal tibiotarsus and the caudally closed ilioischiadic foramen. The actual phylogenetic tree depicts them in a polytomy with both Galloanserae and Neoaves. It has been suggested at times that the "teeth" of pelagornithids were homologous with true teeth on an at least molecular level, being derived from the same programs responsible for the formation of teeth in other dinosaurs. This might have an importance to their actual phylogenetic position. A 2022 paper described Janavis, an ichthyornithine (advanced stem-bird) with a pterygoid bone similar to that of galloanserans. This implies that a galloanseran-like pterygoid is ancestral for crown-group birds as a whole, rather than a derived feature of neognaths. The authors noted that among the groups often regarded as galloanserans based on their pterygoid morphology (pelagornithids, dromornithids and gastornithids), some might instead constitute early-diverging crown-birds outside Galloanserae, or even be outside the avian crown group altogether. Genera and unidentified specimens Due to the fragmented and crushed state of most pseudotooth bird remains, it is not clear whether the roughly one dozen genera that have been named are all valid. Only the beaks are robust and distinctive enough to allow for good taxonomic delimitation, and even these are usually found as broken pieces. For example, Argilliornis and Neptuniavis were recently found to be arm and leg bones, respectively, of Dasornis, which until then was only known from skull bones. Size is generally regarded as reliable marker for generic diversity, but care just be taken to ascertain whether smallish specimens are not from young birds. Tentatively, the following genera are recognized: Protodontopteryx (Early Paleocene of New Zealand) Pseudodontornis (Late Paleocene ?–? Late Oligocene of Charleston, South Carolina, US) – polyphyletic (type species in Palaeochenoides/Pelagornis)? "Odontoptila" (Late Paleocene/Early Eocene of Ouled Abdoun Basin, Morocco) – a nomen nudum; preoccupied Odontopteryx (Late Paleocene/Early Eocene of Ouled Abdoun Basin, Morocco – Middle Eocene of Uzbekistan) – including "Neptuniavis" minor, may include "Pseudodontornis" longidentata, "P." tschulensis and Macrodontopteryx Dasornis (London Clay Early Eocene of Isle of Sheppey, England) – including Argillornis, "Lithornis" emuinus and "Neptuniavis" miranda; may include "Odontopteryx gigas" (a nomen nudum), "Pseudodontornis" longidentata and Gigantornis Macrodontopteryx (London Clay Early Eocene of England) – may include "Pseudodontornis" longidentata and/or belong in Odontopteryx cf. Odontopteryx (Early Eocene of Virginia, US) Gigantornis (Ameki Middle Eocene of Ameki, Nigeria) – may belong in Dasornis cf. Odontopteryx (Middle Eocene of Mexico) Pelagornithidae gen. et sp. indet. (Middle Eocene of Mount Discovery, Antarctica) – same as large Seymour Island specimen/Dasornis/Gigantornis? Pelagornithidae gen. et sp. indet. (Middle Eocene of Etterbeek, Belgium) – Dasornis/Macrodontopteryx? "Aequornis" (Middle Eocene of Kpogamé-Hahotoé, Togo) – a nomen nudum Pelagornithidae gen. et spp. indet. (La Meseta Middle/Late Eocene of Seymour Island, Antarctica) – two species? Same as Mount Discovery specimen/Dasornis/Gigantornis, Odontopteryx? Pelagornithidae gen. et sp. indet. (Late Eocene of France) Pelagornithidae gen. et sp. indet. (Late Eocene of Kazakhstan) – may belong in Zheroia Pelagornithidae gen. et sp. indet. (Eocene of South Shetland Islands, South Atlantic) cf. Dasornis (Late Eocene/Early Oligocene of Oregon, US) – Cyphornis? cf. Macrodontopteryx (Early Oligocene of Hamstead, England) – may belong in Proceriavis Pelagornithidae gen. et sp. indet. (Early Oligocene of Japan) Caspiodontornis (Late Oligocene of Pirəkəşkül, Azerbaijan) – may belong in Guguschia Palaeochenoides (Late Oligocene of South Carolina, US) – may include Pseudodontornis longirostris or belong in Pelagornis Pelagornithidae gen. et sp. indet. (Late Oligocene of South Carolina, US) Pelagornithidae gen. et sp. indet. (Yamaga Late Oligocene of Kitakyushu, Japan) – Osteodontornis? Tympanonesiotes (Late Oligocene or Early Miocene of Cooper River, US) Cyphornis (Early Miocene of Carmanah Point, Vancouver Island, Canada) – may include Osteodontornis Osteodontornis (Early Miocene – Early Pliocene) – may belong in Cyphornis Pelagornis (Early Miocene of Armagnac, France – Early Pleistocene of Ahl al Oughlam, Morocco) – may include Pseudodontornis longirostris, Palaeochenoides Pelagornithidae gen. et spp. indet. (Early? Miocene – Early Pliocene of eastern US) – 2–3 species? Pelagornis? cf. Osteodontornis (Capadare Middle Miocene of Cueva del Zumbador, Venezuela) cf. Osteodontornis/Pelagornis (?Middle/Late Miocene of North Canterbury, New Zealand) cf. Pelagornis (Bahía Inglesa Middle Miocene of Chile – Early Pliocene of Chile and Peru) – 2 species? cf. Osteodontornis (Pisco Middle Miocene –? Early Pliocene of Peru) – 2 species? "Pseudodontornis" stirtoni (Miocene or Pliocene of Motunau Beach, New Zealand) – sometimes Neodontornis Pelagornithidae gen. et sp. indet. (Yushima Early Pliocene of Maesawa, Japan) – Osteodontornis? cf. "Pseudodontornis" stirtoni (Tangahoe Mudstone Middle Pliocene of Hawera New Zealand) Pelagornithidae gen. et sp. indet. (Dainichi Early Pleistocene of Kakegawa, Japan) – Osteodontornis? Pelagornis sp. (Late Pliocene of California, US: Boessenecker and Smith; 2011) Some other Paleogene (and in one case possibly Late Cretaceous) birds, typically taxa known only from the most fragmentary remains, might also be pelagornithids. They are not usually placed here, but the fossils' large size and the known similarities of certain pseudotooth birds' bones to those of other lineages warrant further study. The genera in question are Laornis, Proceriavis, Manu and Protopelicanus.	Biology and health sciences	Prehistoric birds	Animals
7082456	https://en.wikipedia.org/wiki/Astronomical%20Almanac	Astronomical Almanac	The Astronomical Almanac is an almanac published by the United Kingdom Hydrographic Office; it also includes data supplied by many scientists from around the world. On page vii, the listed major contributors to its various Sections are: H.M Nautical Almanac Office, United Kingdom Hydrographic Office; the Nautical Almanac Office, United States Naval Observatory; the Jet Propulsion Laboratory, California Institute of Technology; the IAU Standards Of Fundamental Astronomy (SOFA) initiative; the Institut de Mécanique Céleste et des Calcul des Éphémerides, Paris Observatory; and the Minor Planet Center, Cambridge, Massachusetts. It is considered a worldwide resource for fundamental astronomical data, often being the first publication to incorporate new International Astronomical Union resolutions. The almanac largely contains Solar System ephemerides based on the JPL Solar System integration "DE440" (created June 2020), and catalogs of selected stellar and extragalactic objects. The material appears in sections, each section addressing a specific astronomical category. The book also includes references to the material, explanations, and examples. It used to be available up to one year in advance of its date, however the current 2024 edition became available only one month in advance; in December 2023. The Astronomical Almanac Online was a companion to the printed volume. It was designed to broaden the scope of the publication, not duplicate the data. In addition to ancillary information, the Astronomical Almanac Online extended the printed version by providing data best presented in machine-readable form. The 2024 printed edition of the Almanac states on page iv: "The web companion to The Astronomical Almanac has been withdrawn as of January 2023." Publication contents Section A PHENOMENA: includes information on the seasons, phases of the Moon, configurations of the planets, eclipses, transits of Mercury or Venus, sunrise/set, moonrise/set times, and times for twilight. Preprints of many of these data appear in Astronomical Phenomena, another joint publication by USNO and HMNAO. Section B TIME-SCALES AND COORDINATE SYSTEMS: contains calendar information, relationships between time scales, universal and sidereal times, Earth rotation angle, definitions of the various celestial coordinate systems, frame bias, precession, nutation, obliquity, intermediate system, the position and velocity of the Earth, and coordinates of Polaris. Preprints of many of these data also appear in Astronomical Phenomena. Section C SUN; covers detailed positional information on the Sun, including the ecliptic and equatorial coordinates, physical ephemerides, geocentric rectangular coordinates, times of transit, and the equation of time. Section D MOON: contains detailed positional information on the Moon including phases, mean elements of the orbit and rotation, lengths of mean months, ecliptic and equatorial coordinates, librations, and physical ephemerides. Section E PLANETS: consist of detailed positional information on each of the major planets including osculating orbital elements, heliocentric ecliptic and geocentric equatorial coordinates, and physical ephemerides. Section F NATURAL SATELLITES; covers positional information on the satellites of Mars, Jupiter, Saturn (including the rings), Uranus, Neptune, and Pluto. Section G DWARF PLANETS AND SMALL SOLAR SYSTEM BODIES: includes positional and physical data on selected dwarf planets, positional information on bright minor planets and periodic comets. Section H STARS AND STELLAR SYSTEMS: contains mean places for bright stars, double stars, UBVRI standards, ubvy and H beta standards, spectrophotometric standards, radial velocity standards, variable stars, exoplanet and host stars, bright galaxies, open clusters, globular clusters, ICRF2 radio source positions, radio flux calibrators, x-ray sources, quasars, pulsars, and gamma ray sources. Section J OBSERVATORIES: was a worldwide index of observatory names, locations, MPC codes, and instrumentation in alphabetical order and by country. This section has now been removed as stated in the printed 2024 edition on page J1: "We are presently reserving Section J for possible new contents in future editions of The Astronomical Almanac." An explanation is given on page iv: "Section J: Observatories: This section has been removed as it is significantly out-of-date and it is not clear that a static listing of Observatories is a useful service any longer." Section K TABLES AND DATA: includes Julian dates, selected astronomical constants, relations between time scales, coordinates of the celestial pole, reduction of terrestrial coordinates, interpolations methods, vectors and matrices. Section L NOTES AND REFERENCES: gives notes on the data and references for source material found in the almanac. Section M GLOSSARY: contains terms and definitions for many of the words and phrases, with emphasis on positional astronomy. Publication history The Astronomical Almanac is the direct descendant of the British and American navigational almanacs. The British Nautical Almanac and Astronomical Ephemeris had been published since 1766, and was renamed The Astronomical Ephemeris in 1960. The American Ephemeris and Nautical Almanac had been published since 1852. In 1981 the British and American publications were combined under the title The Astronomical Almanac." Explanatory Supplement to the Astronomical Almanac The Explanatory Supplement to the Astronomical Almanac, currently in its third edition (2013), provides detailed discussion of usage and data reduction methods used by the Astronomical Almanac. It covers its history, significance, sources, methods of computation, and use of the data. Because the Astronomical Almanac prints primarily positional data, this book goes into great detail on techniques to get astronomical positions. Earlier editions of the supplement were published in 1961 and in 1992.	Technology	Astronomical technology	null
7086534	https://en.wikipedia.org/wiki/Kelvin%27s%20circulation%20theorem	Kelvin's circulation theorem	In fluid mechanics, Kelvin's circulation theorem states:In a barotropic, ideal fluid with conservative body forces, the circulation around a closed curve (which encloses the same fluid elements) moving with the fluid remains constant with time. The theorem is named after William Thomson, 1st Baron Kelvin who published it in 1869. Stated mathematically: where is the circulation around a material moving contour as a function of time . The differential operator is a substantial (material) derivative moving with the fluid particles. Stated more simply, this theorem says that if one observes a closed contour at one instant, and follows the contour over time (by following the motion of all of its fluid elements), the circulation over the two locations of this contour remains constant. This theorem does not hold in cases with viscous stresses, nonconservative body forces (for example the Coriolis force) or non-barotropic pressure-density relations. Mathematical proof The circulation around a closed material contour is defined by: where u is the velocity vector, and ds is an element along the closed contour. The governing equation for an inviscid fluid with a conservative body force is where D/Dt is the convective derivative, ρ is the fluid density, p is the pressure and Φ is the potential for the body force. These are the Euler equations with a body force. The condition of barotropicity implies that the density is a function only of the pressure, i.e. . Taking the convective derivative of circulation gives For the first term, we substitute from the governing equation, and then apply Stokes' theorem, thus: The final equality arises since owing to barotropicity. We have also made use of the fact that the curl of any gradient is necessarily 0, or for any function . For the second term, we note that evolution of the material line element is given by Hence The last equality is obtained by applying gradient theorem. Since both terms are zero, we obtain the result Poincaré–Bjerknes circulation theorem A similar principle which conserves a quantity can be obtained for the rotating frame also, known as the Poincaré–Bjerknes theorem, named after Henri Poincaré and Vilhelm Bjerknes, who derived the invariant in 1893 and 1898. The theorem can be applied to a rotating frame which is rotating at a constant angular velocity given by the vector , for the modified circulation Here is the position of the area of fluid. From Stokes' theorem, this is: The vorticity of a velocity field in fluid dynamics is defined by: Then:	Physical sciences	Fluid mechanics	Physics
2953081	https://en.wikipedia.org/wiki/Limousin%20cattle	Limousin cattle	The Limousin, , is a French breed of beef cattle from the Limousin and Marche regions of France. It was formerly used mainly as a draught animal, but in modern times is reared for beef. A herd-book was established in France in 1886. With the mechanisation of agriculture in the twentieth century, numbers declined. In the 1960s there were still more than 250 000 head, but the future of the breed was not clear; it was proposed that it be merged with the other blonde draught breeds of south-western France – the Blonde des Pyrénées, the Blonde de Quercy and the Garonnaise – to form the new Blonde d'Aquitaine. Instead, a breeders' association was formed; new importance was given to extensive management, to performance recording and to exports. In the twenty-first century the Limousin is the second-most numerous beef breed in France after the Charolais. It is a world breed, raised in about eighty countries round the world, many of which have breed associations. History The Limousin originates in the Limousin, the area surrounding Limoges on the western flank of the Massif Central. It was originally a robust draught animal, used for agricultural work. As elsewhere, oxen at the end of their working lives were fattened and sent to slaughter, at times in major cities such as Bordeaux or Paris. In 1791, Jacques-Joseph Saint-Martin, an agronomist from Limoges, acknowledged the importance of Limousin cattle in the markets of cities such as Paris, Lyon, and Toulouse. Limousin cattle actually came from the departments of Charente, Dordogne, Haute-Vienne, Vienne, Lot, Corrèze, and Creuse. The market for Limousin cattle declined slightly in the early 19th century, but livestock still remained a major activity in the region. A large variation in the agricultural systems was operating in the Limousin region, defined by three types of district. These were productive, grain-producing areas, called d'engrais, undeveloped, marginal, predominantly forested land called forestiers, and developing land called d'élèves. Cattle, in particular cows, were used extensively for all types of agricultural work. At the beginning of the 19th century, the Limousin region was characterised by the mediocrity of its animals. Texier-Olivier Louis, prefect of the Haute-Vienne, observed that Limousin cattle weighed 300 to 350 kg and measured 1.5 m at the withers. The defect was considered to be attributable to poor genetics, nutrition and breeding practices. In competitions, Limousin cattle were among the worst performers. The breed was considered to be a working breed, but poorly shaped and with poor milk. To improve the breed, some Limousin breeders tried to cross their animals with Agenais, Norman or Charolais cattle, which were better shaped. The Limousin breed was also not immune to the wave of Anglomania affecting France in the middle of the 19th century. Some wealthy farmers maintained Durham cattle, which were prized by the elite of the time. However, this practice was criticised by the agricultural society of Limoges. The society encouraged farmers to continue selecting animals that were most consistent with the characteristics of the Limousin breed, which was perfectly adapted to the region's environment, rather than trying to adapt other breeds. Furthermore, the vast majority of Limousin breeders could not afford to raise livestock in addition to their working animals, as was the case on larger properties that practiced crossing with Durham cattle. Finally, the marginalisation of English animals in competitions from the late 1860s reinforced the case to improve the breed by itself. A herd-book was started in 1886. At the beginning of the 19th century, a bonus was introduced to reward farmers who retained their best bulls, though they were not productive. The increase in weight of the animals began with improved grassland. The second half of the 19th century had the arrival of fertilizers and improved pastures such as clover and ryegrass, which not only improved the productivity of existing fields, but also transformed the moorland pasture. Vineyards affected by phylloxera were also being returned to pasture. The results were not immediate. In 1862, cattle sold at La Souterraine weighed about 600 kg. The decline of Anglomania in favour of economic pragmatism, and the criticism and fall of the aristocracy aided the development of Limousin cattle. The crowning moment was the honour received by the bull Achilles Caillaud to open the competition in Paris for all breeds in 1886 (the year the Limousin Herd Book was created), and the grand champion prize of all breeds won three years later by Charles Léobardy for his team. The First World War slowed down the growth of the Limousin breed, which carried through the interwar years despite a reorganisation of the herd book in 1923. Herd book registrations grew slowly, from 600,000 animals in 1890 to 800,000 in 1940. The Limousin breed almost disappeared when the French government planned to combine it with the Garonne, Quercy, and Blonde des Pyrenées breeds, during the formation of the Blonde d'Aquitaine breed in 1962. All of these cattle belonged to the "blond and red" branch of cattle. Limousin breeders fiercely opposed the merger and the Limousin breed was retained. The Limousin breed resumed its growth in the 1960s. The size of the French Limousin herd has increased sharply in recent years, with a 50% increase in numbers in France in 15 years. Today, it is the second-most numerous French beef breed, behind Charolais and ahead of Blonde d'Aquitaine. In 2004, of about 900,000 Limousin cows, 63,000 were recorded in the herd book. At that time, 20,000 bulls were used for breeding, 10% through artificial insemination, and 1,600 were recorded in the Herd Book. France's Limousin herd grew by 2.6% in 2014 to 2.69 million head as of 31 Dec., including 1.09 million cows. DNA studies have identified close genetic relationships between Limousin cattle and other south-west European breeds. One study reported a possible common origin or recent gene flow between the Limousin and Charolais cattle breeds. whereas other studies indicated that a closer genetic relationship exists between Limousin, Gasconne, Aubrac, Bazadais, Salers, and Blonde d'Aquitaine cattle. One historian reported that the Limousin breed's origins can be traced to the blonde Garonne breed in the fifth century AD. The Garonne breed from the south-west of France was merged into the Blonde d'Aquitaine breed in 1962. The grey Gasconne breed with which Limousin cattle have a close genetic relationship is also reported to have arrived in the south-west of France with the Visigoths also around the fifth century AD. The Limousin belongs to the blond group of European cattle, in a sub-group which also includes the Alpha 16, the Maraîchine, the Nantaise and the Parthenaise. The herd-book Significance The purest form of Limousins have ancestors that can all be traced to "Full French" entries in the herd-book (known in France as Le livre généalogique). These Limousins are known by different names. In the US, and Canada they are known as Fullbloods, in Australia and New Zealand as French Pure, and in European countries such as Britain as purebred or simply Limousin. In France, two Full French Herd Book classes exist, namely Pureblood (pur sang in French, also translated to Fullblood) and Pure Bred (race pure in French). The Full French Pure Bred Herd Book class, as with all European Union (EU) member countries' herd books, is controlled by EU legislation. Full French is a term used by the French Limousin breeders' association (known in France as Herd Book Limousin, abbreviated to HBL) to describe cattle that comply with: Bred by French active member-associates of the HBL: The strictly enforced rules of the HBL require breeders to conduct on-farm performance testing of their animals and to have selected animals independently tested by approved official bodies. Independently inspected and certified to be Full French according to the Breed Standard. Cattle excluded from Full French certification include those imported into France, cattle that are polled (in French sans corne), and cattle that have undesirable double muscling genes (in French gene culard) inherited from non-Limousin base animals. A less pure form of Limousin is bred up (also known as graded up) from a base animal over a defined number of generations. A parent of each generation's progeny must be registered as a Limousin in the respective country's herd book. In the US, Canada, Australia, and New Zealand, a graded up Limousin, after three generations for females and four generations for males, is known (confusingly with the legal European definition) as purebred, which is then eligible for recording in the respective countries' herd books alongside Fullblood and French Pure Limousins. Unlike the US, Canada, Australia, and New Zealand, which allow both purebred and Full French bulls and dams to be used for grading up, in Britain, grading up can only occur using Full French bulls. British graded up females when they reach fourth generation from a non-Limousin base cow can then be registered as Limousins in the British Limousin Pedigree Register. The British Limousin Pedigree Register is separate from the British Limousin Herd Book, which is reserved for animals that have complete Full French ancestry. Evolution of Herd Book Prior to July 2007 The herd-book was destroyed in the Second World War. When the French Government decided to merge the Limousin breed into the new Blonde d'Aquitaine breed in the 1960s, which was fiercely opposed by French Limousin breeders, the impetus to re-establish the herd-book was provided. Inspectors were appointed to identify "true to type" Limousins from the Limousin region. These were admitted to the new herd-book as foundation animals (in French titre initial, abbreviated to T.I.). Following its re-establishment, the Herd Book was opened from time to time for the admission of T.I. females that on inspection satisfied the Full French breed standard. These animals were identified by the letters T.I. placed after their name. The process of admitting new T.I. animals to the Herd Book continued until June 2008. The Limousins recorded in the herd-book were known as Pureblood (literal translation of the French pur sang). The is normally the name given to English thoroughbred horses, although in the context of Limousins the English translation Fullblood is commonly used. July 2007 to June 2008 EU legislation, pressure from French breeders of polled stock, and other developments, including requirements of European Limousin associations (the 11 countries of EUROLIM), contributed to a restructuring of the herd-book that commenced in July 2007. During the period July 2007 to June 2008, the herd-book comprised a main section (section principale in French) divided into the original Pureblood (pur sang) class and a newly created Purebred (race pure) class. The Purebred class was added to enable the recording of polled animals, those that carried a double-muscling gene (muscle hypertrophy abbreviated to MH, or gene culard in French), and those that did not comply fully with the French Breed standard. After June 2008 EU legislation allowed a supplementary section (section annexe in French) to be used to introduce genetics into existing breeds from other breeds in a grading up process aimed at "progressive improvement". According to the legislation, only females whose mother and maternal grandmother entered in a supplementary section, and whose father and two grandfathers are entered in the main section, can be regarded as purebred and entered in the main section of a herd book. Although this appears to be a simple two-stage grading up process, base females that start a new grading up line were also required by EU legislation to "be judged to conform to the breed standard". Since 2007, EU legislation allowed base animals to be males but infusion of their genetics into the main section is still only possible through their female progeny. The restructured French Herd Book is described as having a third section called certified purebred (race pure certifié in French) intermediate between the first two for recording animals that do not comply with the breed standard (for example incorrectly coloured hair in certain places), have double muscling genes, or are polled. Limousins imported into France that comply with Council Directive 2009/157 are also recorded in sub-class 2 (sous-classe 2 in French) of the certified purebred class because they do not comply with the French HBL requirement of being Full French. Base animals selected for the two-stage grading up process to any EU herd book purebred class require a minimum of 3/4 Limousin content. Graded up females using the two-stage process then become eligible for entry into the main section of all EU herd book purebred classes as initial registration (or T.I.) Limousins when they reach 15/16 Limousin content. They are then legally identified as Limousin (i.e. 100% Limousin) – the breed code 34 often substitutes for the word Limousin in French discussions and reports on cattle breeding. Only the Certified Purebred sub-class 2 and Registered Purebred class are identified as Limousin in France because cattle of non-Limousin origin had been introduced into the supplementary section of the Pureblood class. The growth and spread of the Limousin breed in France since the early 1980s meant that a past practice of selecting a base female on appearance alone was no longer a guarantee of its breed origin because of the potential for "crossing absorption". Base females inadvertently assessed as Limousin and recorded as T.I. animals in the main section of the Pureblood class included Parthenais and Charolais, which were presumably the source of the double-muscling genes found in the French Limousin Pureblood population. In 2008, the double muscling gene had been found in 3% of active bulls in France. Currently, only cattle recorded in the French Limousin Certified Purebred sub-class 2 and Registered Purebred class satisfy the requirements of EU law on herd books and can be transferred as Limousins, including indirectly through their genetics (for example semen and embryos), to other EU countries and recorded in the respective herd books. Outside of the EU, the rules and regulations of Limousin breed associations do not yet differentiate between the French Pureblood and Purebred classes, with the Pureblood class remaining the origin of, and standard for, the purest form of Limousin. Since the 1960s and until 2008, the French Pureblood class defined the standard against which Limousins throughout the world were measured. Although it would seem appropriate to preserve the integrity of the herd-book as the international Limousin breed standard by preventing the EU-mandated entry of animals that do not meet the Full French standard, restrictions to such entries remain forbidden under EU law and subsequent interpretations by the European Court of Justice. Immediately prior to the restructuring of the herd-book in 2008, French breeders had two months to nominate the class (Pureblood or Purebred) in which they wanted their cattle to be registered. Pureblood cattle have higher business value than Purebreds because they are preferred by French breeders. Also, Full French cattle benefit under a branded meat marketing label known as Blason Prestige Label Rouge, which has around 800 outlets in France. Future EU legislation on dehorning practices is expected to place greater pressure to breed polled cattle, which is being pursued in France by the GIE Polled Excellence group. Because no French Limousins had been identified with the polled gene, a breeding programme commenced in 2005 with polled Canadian bulls of phenotype closest to the French breed standard. French breeders of polled Limousins claimed that the breed standard that prevented their animals being recorded in the French Herd Book provided an unfair export advantage to foreign countries that do not have similar restrictions. The French recording ban was removed in July 2007 with the introduction of the Purebred class, but so far no polled Limousin have yet been accepted and registered as Full French. Characteristics Limousin breed standard The French Limousin breed standard is described in Article 1 of Title I of the Rules of Procedure of the French Limousin herd-book, 1 August 1991: The Limousin is a large framed breed of beef cattle with a bright wheat-coloured coat, not too dark, a little lighter on the belly, the rear of the thighs, between the legs, on the anus, around the testicles or udder, and the tail tip. Absence of any spots or pigmentation, with pink mucous membranes. Short head, broad forehead and muzzle, lighter area around the eyes and muzzle, fine horns curved forward and slightly raised at the tip (if present). Short neck. Chest broad and rounded. Side round. Pelvis wide, especially at the pin bones, not too inclined. Bones of lower back and hips slightly protruding. Forequarter well-muscled, wide above and very muscular. Hindquarters thick, deep and rounded. Horns and hooves lighter coloured. Correct limbs. Fine supple hide. Characteristics considered unacceptable in the French breed standard: Any pigmentation or black spots on muzzle, black or white hairs anywhere on the coat, particularly in the ears, at the end of the tail and around the muzzle. White hairs anywhere. An eliminating count of less than five for any of the different breed standard points. A difficult or vicious disposition (tranquilising is forbidden). Any obvious physical malformation. The French Limousin breed standard is applied by international breed associations in slightly different forms. These range from mandatory compliance before an animal can be recorded in a country's herd book (mainly European countries) to voluntary application in others. For example, in Belgium, application of its breed standard mirrors in most detail the French use, and in the UK, compliance with its version of the Limousin breed standard is required by the UK breed association's bye-laws. The USA, Canadian, Australian and New Zealand breed associations do not define a breed standard in their regulations, and application of any standard is voluntary. The only requirement for registration as a Fullblood in both North American herd book registers is that ancestors should have "full French ancestry", or trace directly to the "Herd Book Limousin in France". In Australia and New Zealand the French Pure herd book classification requires that animals carry "100% Pure French genetics". USA, Australian and New Zealand breed association regulations also allow graded up animals to be registered in their herd books as purebreds without a requirement to comply with a minimum French Limousin content. Grading up using these purebreds over base or lower grade animals has resulted in the gradual reduction in the French Limousin content of some purebreds, and an observable divergence from the French breed standard. The Canadian breed association by regulation prevents loss of French Limousin content from its registered purebreds by requiring that they "contain 90% or more Limousin blood". In the US, Canada, Australia and New Zealand, breed development and selection is influenced principally by performance recording and genetic selection. French performance recording and genetic prediction The breed standard in France is applied in parallel to an intensively applied system of selection, performance recording and genetic prediction that was established gradually across the country commencing in the 1980s. The system appears to be similar to that used in Denmark. All females recorded in the French Herd Book are controlled under this system, which focuses mainly on maternal qualities derived from measurements of calving ease, and growth and structure of calves. Females that achieve the best indexes for particular traits are then examined in detail by HBL technicians who assign grades based on morphology. The best females receive the qualification Reproductive Recognised (in French Reproductrice Reconnue, abbreviated to RR), which is awarded to the top 10%, or Reproductive Recommended (in French Reproductrice Recommandée, abbreviated to RRE) awarded to the top 1%. The qualifications aid the identification of superior animals. For males, selection of the best breeders is more complex. The first step is weaning, when the morphology of calves and the known qualities of their parents are used to make an initial selection of animals that receive the qualification Reproductive Hope (in French Reproducteur Espoir, abbreviated to Espoir). Annually in excess of about 700 bull calves are then selected to enter the national evaluation station at , close to Limoges, just after weaning, when they are about seven months old. At Lanaud the animals are grouped together to compare their performance under identical feeding and environmental conditions to the age of 13–14 months. The differences observed between the animals are then related principally to their genetics, which is of interest to breeders because this is what is transmitted to a bull's progeny. After completing evaluation at Lanaud, half of the young bulls are awarded the qualification Reproductive young (in French Reproducteur jeune, abbreviated to RJ) by the HBL. Most of these bulls are intended for natural service, and the best are subjected to evaluations of their progeny. In the same way as for females, the best bulls receive the qualification "Reproductive Recognised" (RR), awarded to the top 10%, or "Reproductive Recommended" (RRE), awarded to the top 1%. In parallel with the Lanaud evaluation station there are three local stations at La Souterraine in the Creuse department of the Limousin region, Saint-Jal in Corrèze, also in the Limousin region, and Naucelle in Aveyron in the south of France. The local stations provide commercial beef producers in their region with bulls of high production potential for use by commercial farmers whose herds are not necessarily registered in the French herd book. The best bulls go to artificial insemination (AI) cooperatives where semen is taken. AI allows the wide distribution of a bulls' genetics to the benefit of more farmers. However, in order to guarantee their genetic qualities, the bulls are subject to a strict selection scheme to increase the accuracy of the different genetic indexes. The best bulls identified at Lanaud are sent to another test station at Naves in Corrèze. Here they are tested more accurately and evaluated for feed conversion efficiency, growth and muscle development. Progeny of the top 10 bulls out of this testing, and the best natural service bulls, are then evaluated. Cows are inseminated to produce 60 to 80 calves per bull tested, which are in turn evaluated. Male progeny go to a station in Pépieux in the south of France, where they are fed a ration of corn silage before being slaughtered at the age of 16 months. In addition to evaluations of growth and conformation in the live animals, carcases, including fat composition, are evaluated. The best bulls identified in progeny testing are formally given the award Young Beef Cattle (in French Viande Jeunes Bovins, abbreviated to JB). Female progeny go to a test station in Moussour in Corrèze, where they are inseminated with the same bulls and calve at two years in confinement before being put out to pasture with their calves. The test station evaluates weight, growth, morphology, fertility, calving ability and milking ability in order to assess their maternal qualities. The best bulls following the tests on their daughters are identified as Maternal Qualities (in French Qualités Maternelles, abbreviated to QM). The qualifications RR and RRE are recorded with an animal's description in sales' catalogues and other promotional literature. As a further aid to purchasers of French Limousin genetics, additional qualifications provide a guide to the greatest likely production benefit based on an animal's genetics estimated from on-farm progeny testing. The qualifications are aligned with French market specifications for Limousin beef: VS – awarded to sires recognised or recommended for weaner production. VB – awarded to AI sires recognised or recommended for vealer production. JB – awarded to AI sires recognised or recommended for general beef production. QM – awarded to AI sires recognised or recommended for breeding stock production. M – awarded to AI sires recognised or recommended for both general beef and breeding stock production. P – awarded to females recognised or recommended for the production of early developing progeny. T – awarded to females recognised or recommended for the production of late developing progeny. Genetic basis for muscling in Limousin cattle The Limousin cattle breed has been popular in France for more than two centuries because of its meat qualities and the breed's production efficiency. Since the early 1990s scientists have quantified these breed characteristics in comparisons with other breeds, and identified a natural variant of the myostatin gene found in Limousins which has a significant influence on them. The myostatin gene is found in all mammals and influences the production of a protein that controls muscle development. Variants of the gene produce proteins that are less effective at controlling muscle development, which results in increased muscle mass. Limousin muscling is intermediate to that of British cattle breeds such as Aberdeen Angus, Hereford, and Shorthorn and the extreme double muscling found in the European Belgian Blue and Piedmontese breeds. Studies of double-muscled cattle identified natural mutations of the myostatin gene which produce inactivated proteins that are unable to control muscle development. In Belgian Blue and Piedmontese cattle this causes an increase in muscle mass of 20–25%. Subsequent studies identified a less extreme myostatin mutation known as F94L associated with Limousins. The resulting partially active protein results in Limousins having intermediate muscle development, which avoids the extreme muscling and associated disadvantages of double muscled cattle. A Limousin/Jersey backcross study conducted in Australia and New Zealand to investigate the effects of the F94L myostatin variant concluded that the mutation had no significant effect on birth-weight and growth traits. Averaged over all backcross calves in the trial (total of 766), animals homozygous for the mutation had approximately 6% heavier carcases than animals without the mutation, 15% larger eye muscle (also known as rib eye) area, 13% heavier silverside weight, and 13% heavier total meat weight. Increased meat weight and size was accompanied by a 15% reduction in intra-muscular fat and 25% reduction in total fat weight. No other significant effects were observed. A second backcross study conducted in Japan of Limousin and Japanese Black breeds identified similar changes to meat and fat quantities in cattle homozygous for the F94L mutation. Although the Australian/New Zealand study found that the F94L mutation was partially to significantly recessive in most traits, meaning cattle heterozygous for the mutation express less to significantly less than half of the effects noted for homozygous cattle, the Japanese study found that the meat and fat quantities in cattle heterozygous for the mutation were about midway between the two extremes. Distribution of F94L myostatin variants in Limousin cattle A number of international breed associations have been testing the F94L status of cattle registered in their herd books. The absence of F94L genes in some tested cattle might be a result of a sire or dam ancestor that had double muscling (MH) genes, or more likely that the myostatin gene was the normal or wild type variant found most commonly in beef cattle. In the latter case, loss of the F94L variant will arise when grading up to purebred when base animals are not Limousins. Europe Of the 14 Limousins tested during the research that led to the discovery of the F94L variant, 12 cattle were homozygous for the variant and two were heterozygous. The second myostatin genes in both heterozygous cases were each different myostatin MH variants of types normally found in Belgian Blue and Charolais cattle. Britain British test results of sale bulls in February 2010 indicated that of 142 animals tested, just under 90% were homozygous for the F94L mutation, about 8.5% were heterozygous, and 1.5% did not have the mutation. USA Test results of approximately 1,100 cattle recorded in the North American Limousin Foundation herd book show the following distributions for three classes of animal. About 94.4% of Fullbloods, 62.3% of purebreds and 5.3% of Lim-Flex were homozygous for the F94L mutation. Australia and New Zealand Test results of 1028 cattle recorded in the Australian and New Zealand herd book indicate that 96.7% of Fullbloods (known in Australia as French Pure), 88.0% of purebreds, and 33.3% of a limited sample of Lim-Flex were homozygous for the F94L mutation. Implications for cattle heterozygous for the F94L variant Cattle heterozygous for the F94L myostatin mutation have a 50% probability of passing the mutation to their progeny. Because the mutation has greatest effect on carcase traits, only 50% of progeny of a heterozygous parent will inherit increased muscling associated with the mutation. Furthermore, best linear unbiased prediction (BLUP) techniques used to estimate the genetic merit of stud cattle (for example, estimated breeding values (EBVs) and expected progeny differences (EPDs)) will be incorrect because they assume that no dominant genes contribute to modelled traits. Inconsistent inheritance of myostatin mutations (for example, F94L in Limousins, nt821 in Angus, and Q204X in Charolais) by progeny is expected to result in possible BLUP prediction errors for EBVs and EPDs equalling or exceeding worst case standard errors of prediction. For example, average rib eye area for Limousins in US Meat Animal Research Center (USMARC) trials during the 1980s and early 1990s is reported to be 12.3in2, and the reported possible difference in rib eye area in progeny arising from inheritance of either two F94L mutations or two normal myostatin genes from heterozygous parents is estimated to be 1.8in2 (12.3in2 x 15%). This difference, which is unpredictable without DNA testing, is nearly four times the possible change value for a 0% BIF accuracy, reported to be 0.46in2 for the rib eye EPD. When one parent is heterozygous for the mutation, and the other homozygous for the mutation or the normal form of the myostatin gene, the expected average difference in rib eye area of progeny will be about 0.9in2 (12.3in2 x 7.5%), depending on whether the mutation or normal form of the gene is inherited from the heterozygous parent. In this case, the unpredictable variation in rib eye area represents about twice the possible change value for a 0% BIF accuracy. Standard errors of prediction, also known as accuracy or possible change value in the context of EBV and EPD predictions, are dependent on the quality of information used to predict an animal's EBV or EPD for a given trait. Errors in estimating genetic merit are being addressed in research programmes that aim to supplement phenotypic data extensively used in current BLUP predictions with genotypic data. Comparisons with other breeds A USMARC long-term multi-breed study of Limousins, three British (Red Poll, Hereford, Aberdeen Angus) and five other continental European (Braunvieh, Pinzgauer, Gelbvieh, Simmental, Charolais) cattle breeds reported that Limousin cattle were the most efficient and fastest of all breeds at converting feed into saleable meat even though Limousin's live weight growth was the slowest. This arose because saleable meat yield expressed as percentage of live weight was significantly higher in Limousins than in most other cattle breeds. Saleable meat yield was an average 34.9% of live weight for the three British cattle breeds, compared with 40.4% for the five other continental European breeds, and 46.0% for Limousins, for two market end points of 225 kg saleable meat at 8mm fat trim, and 210 kg saleable meat at 0mm fat trim. Live weight gain for the Limousins averaged 1.27 kg/day, compared with an average 1.29 kg/day for the British breeds and 1.38 kg/day for the other continental European cattle. Limousin saleable meat gain averaged 585g/day, with the balance being 687g/day of low value or waste product, such as bone, trimmed fat, internal organs, and hide. The British breeds produced significantly less saleable meat (average 451g/day) and significantly more low value product (841g/day), while consuming about twice the feed of the Limousins from entry to the trial (weaning) to the market end point (slaughter). The other continental European breeds produced on average less saleable meat (556g/day) and more low cost product (819g/day) while consuming about 25% more feed than the Limousins. Although the Simmental and Charolais produced marginally more saleable meat (590g/day) than Limousins, they produced significantly more low cost product (847g/day) and consumed 18% more feed. For a market end point of 333 kg carcase weight, the Limousin carcases in the USMARC study were estimated to be on average 63.5% of live weight, compared with an average 59.7% (range 58.6% – 60.4%) for the eight other breeds. Similar figures for Limousin meat yield are reported in other sources. The USMARC study indicated that Limousins were significantly the slowest of all breeds to achieve market end points of two measures of marbling score (70 to 160 days longer than the British breeds, and 65 to 70 days longer than the other continental European breeds) while feed conversion efficiency based on live weight gain was marginally poorer (12% less than the British breeds and comparable with the other continental European breeds). When feed conversion efficiency is adjusted to weight of saleable meat divided by feed consumed, Limousin feed conversion efficiency then exceeds both British and continental European breeds by 10–25%. The USMARC study also indicated that Limousins were very significantly the slowest of all breeds to achieve market specifications of three measures of rib eye fat (300 to 400 days longer than the British breeds, and 170 to 220 days longer than the other continental European breeds) while feed conversion efficiency based on live weight gain was poorer (25–30% less than the British breeds and 12–16% less than the other continental European breeds). When corrected to saleable meat divided by feed consumed, feed conversion efficiency was similar or marginally better for Limousins. At these end points, Limousins finished at markedly heavier live weights (up to 490 kg heavier than the British breeds, and 190 kg heavier than the other continental European breeds). The latest USMARC study of Limousins, two of the British breeds and three of the continental European breeds from the original study, reported similar saleable meat yields/live weight for the British breeds (average 36.3%, compared with the earlier 34.9%) and other continental European breeds (average 38.7%, compared with 40.4%), but a significant reduction for Limousin (39.4% compared with 46.0%). However, feed conversion to saleable meat for Limousins for the six reported market end points still exceeded the average of the other two breed groups by up to one-fifth. Live weight and daily live weight gain are the simplest and most common of all traits to be measured and reported, which continues to mask Limousin's superior saleable meat production efficiency. Breed differences are expected to have reduced since the USMARC studies in the 1980s and 1990s because of the wide-scale introduction and use of performance recording and genetic improvement programmes. The reduction in yield reported for Limousins is possibly a result of the loss of French Limousin content and of F94L myostatin mutations from the US purebred population, which would be an expected result of the purebred grading up process practiced there. Earlier USMARC studies are reported to have evaluated high content pure French Limousins. Breed comparison studies of performance-tested bulls report Limousin's more efficient conversion of feed into live weight and confirm the breed's slower live weight gain when compared with other beef cattle breeds. Limousins generally have lower levels of intra-muscular fat (marbling) and subcutaneous fat when compared with British breed cattle grown in similar conditions. Marbling, together with tenderness and flavour, has been associated with eating quality in some countries, and attracts a higher quality grading with associated premiums, although the link between marbling and palatability is not universally supported. In some countries, Limousin's lean, tender meat is considered to be of outstanding quality and marketed through specialist outlets. Beef producers targeting the higher marbling specifications of some markets, but who have concerns over the poorer feed conversion efficiency and yield associated with higher marbling British breed cattle, use Limousin sires over British breed cows, or British breed sires over Limousin cows, in crossbreeding programmes that aim to achieve a balance between the different and conflicting production demands. Crossbreeding with Limousins Crossbreeding increases production efficiency because of hybrid vigour, and allows complementary traits of parents to be combined to produce progeny better suited to different environments or markets. Crossbreeding through the use of Limousin terminal sires in purebred British breed cow herds allows the complementary traits of higher marbling and fat cover provided by the British breed cows, and required or preferred by some markets, to be combined with the higher yield and feed conversion efficiency of Limousin sires. Crossbred cows produce up to, and in some cases in excess of, 20% more weaned calf weight as a result of increased reproductive performance and maternal ability. Crossbred cow longevity is also increased by up to two years when compared with straightbred cows. However, the benefits of hybrid vigour in a crossbred cow decline in subsequent generations if progeny are mated to cattle of parentage similar to the cow, and increase if a new breed is introduced. Although studies acknowledge that the major production benefits of hybrid vigour occur in crossbred cow herds, the main use of Limousins outside of Europe continues to be as terminal sires in purebred British breed cow herds. Genetic basis for crossbreeding Progeny of two parents of different breeds are termed F1 hybrids, F1 crosses or first crosses. F1 hybrids generally have an extremely uniform phenotype and benefit from hybrid vigour. These advantages are observed in the breeding of a wide variety of animals and plants, and arise because progeny inherit one of each paired gene from each parent. When both parents are homozygous for different variants of genes (known as alleles), which is likely to be the case when a breed has been developed and selected over several generations, progeny will inherit both gene variants present in the parents. The F1 hybrid progeny will then be heterozygous for each gene variant, which in turn increases the likelihood that the genes will code for an optimal protein or enzyme. This is the genetic basis of hybrid vigour. While many gene variants have effects that are of little consequence to beef production, a few, such as the myostatin variants found in different cattle breeds, have a major effect. Loss of hybrid vigour occurs and phenotype varies greatly in subsequent generations if F1 hybrids are interbred or backcrossed with animals genetically similar to the F1 parent. Interbred F1 hybrids produce progeny that can be either heterozygous for each gene variant, homozygous for one gene variant, or homozygous for the other gene variant. When one of the variants has a large effect on a trait, for example the effect of myostatin variants on muscularity, larger phenotypic variation will occur among the progeny. Backcross progeny have less phenotypic variation and comprise animals that are either heterozygous for each gene variant or homozygous for the variant found in the original F1 backcross parent. A third form of progeny arises when F1 hybrids are bred with animals genetically dissimilar to their parents. If heterozygosity is maintained or increased as a result, hybrid vigour and other production benefits occurring in the F1 generation will be maintained or increased. Maintenance of heterozygosity is the key to maintaining the highest levels of hybrid vigour. This requires complex breeding programmes and high levels of management. Simplified crossbreeding programmes have been developed using hybrid or composite bulls, which was the motivation behind the development of Lim-Flex hybrids. The two major Limousin hybrids are Brahmousin (a cross between Brahman and Limousin cattle) and Lim-Flex (a cross between Angus and Limousin cattle), which were both developed before the significance of the F94L myostatin variant had been quantified. When Limousins homozygous for the F94L myostatin mutation are used in crossbreeding, only one of the mutations will be inherited (that is, progeny will be heterozygous for the mutation), and a high level of phenotypic uniformity and hybrid vigour would be expected in the progeny. However, breeding using heterozygous animals as parents, which could include purebred Limousins of low percentage Full French content, and Lim-Flex and Brahmousin hybrids that have not been bred to a uniform (homozygous) standard over several generations, would produce progeny with inconsistent carcase characteristics and production value depending upon whether or not the F94L mutation was inherited. The use of Lim-Flex and Brahmousin sires over a third breed of cow would benefit most from increased hybrid vigour, which should minimise any reduction in carcase value arising from the loss of the F94L mutation. According to research into the effects of the F94L mutation, live weights of progeny are unaffected by random inheritance of the mutation. Brahmousin Brahmousin cattle are a hybrid purebred breed of Brahman and Limousin first created in the US in the late 1970s. The goal was to blend the best of the Limousin and Brahman traits to create a breed that has reproductive efficiency, mothering ability, good muscling and growth traits, and adaptability to varying environmental conditions. Brahmousin are now bred in the US, Indonesia, El Salvador, and Australia. The first Brahmousin cattle were produced from a multiple embryo transfer from a French-imported Limousin dam. The resulting progeny were then crossed with Brahman cattle to achieve an F1 hybrid. Further crosses over a broader base led to the production of the 5/8 Limousin – 3/8 Brahman Brahmousin purebred, a mix which has been found to be the most widely accepted and most useful for the majority of the US. The American Brahmousin Council allows animals that are not purebred to be recorded as percentage animals as long as they are at least one-quarter Limousin and one-quarter Brahman. To be recorded as a purebred Brahmousin, the animal must then be sired by a registered purebred or fullblood Limousin bull, registered Brahman bull, or a registered purebred Brahmousin bull. In Australia, Brahmousin are between one-quarter and three-quarters of the parent breeds with the objective of combining the muscle growth and meat quality of Limousins with the heat and parasite resistance, fast growth, and good mothering ability of the Brahman. Brahmousin is formally recognised as a cattle breed in Australia. Lim-Flex Unlike the Brahmousin, Lim-Flex does not have purebred breed status in any participating countries, which includes the US, Australia, New Zealand, and Canada. The need for the Lim-Flex hybrid arose in 2000 out of a perceived need by North American commercial cattle breeders for hybrid bulls that would assist in achieving end-product targets. Lim-Flex is a registered certification mark awarded to Limousin:Angus crossbred or hybrid cattle in the US with content between 25% and 75% Limousin pedigree blood, and between 25% and 75% of either Angus or Red Angus pedigree blood, with a maximum allowable 1/8th of unknown or other breed. Lim-Flex provide genetic options ranging from high content fullblood and purebred Limousin with high levels of muscle and efficiency, to blended options with higher marbling and maternal characteristics associated with Angus cattle, to meet the needs of crossbreeding programmes. The Lim-Flex certification mark has been adopted in Australia and New Zealand, where "commercial Lim-Flex must be 25 to 75 percent Limousin and 25 to 75 percent Angus or Red Angus", and in Canada, where they "must be 37.5 to 75 percent Limousin and 25 to 62.5 percent Angus or Red Angus, with a maximum allowance of another breed or unknown breed composition of 12.5 percent (1/8th)". Appearance Most Limousin cattle's coloration varies from light wheat to darker golden-red. Other coloration, mainly black, has been developed through cross-breeding and grading up from other breeds of cattle. In addition to altering natural coloration, other traits, such as polled (a genetic lack of horns), have been introduced through crossbreeding. Angus cattle have been the dominant source of black and polled genes now present in graded-up, high-content Limousins found throughout the world. Temperament Since the mid-1990s, Limousin breeders have worked extensively on improving disposition issues in Europe, North America, Australia and New Zealand. This has been aided by the high heritability of temperament and by the development of genetic measures of docility (among many other traits) predicted from field measurements and subsequent analysis using BLUP techniques to produce docility EBVs and EPDs. Significant improvement has been recorded in the temperament of Limousins, as can be observed from trait trends published by various international breed associations. Distribution outside France Initial exports Following the creation of the French Limousin Herd-Book in 1886, Limousins were exported to Brazil (1886), New Caledonia (1902), Uruguay (1910), Madagascar (1922), Argentina (1924), and Portugal (1929). However, the only herd that became established outside France during this period was in New Caledonia, which continued to import Limousins. It was not until after the reform of Limousin breeding in France in 1962 that significant numbers were exported around the world. Limousins were reintroduced in Argentina (1966) and Brazil (1978), and imported to other European countries such as Spain (1965), Italy (1968), the Netherlands (1969), Denmark (1970), and the United Kingdom (1971). Their introduction to the United Kingdom provided opportunities for Australia and New Zealand to import semen in 1972. Soon after, New Zealand allowed the importation of Limousins from both Ireland and the United Kingdom, and the first Full French cattle were imported to Australia from New Zealand in 1975. An essential step in the global spread of the Limousin breed was made in 1968 with the importation of the first bull, Prince Pompadour, to Canada. The semen of this bull was made available to the United States in 1969. During the early 1970s, imports of animals to North America started to grow strongly. Today, the North American Limousin Foundation is the largest global Limousin breeders' association. Current situation Limousins ability to adapt to different environments contributed greatly to the breed's current success outside France. In most cases, Limousin bulls or their semen are now imported to improve locally bred Limousins. Today, the breed is present in about 70 countries around the world, and in all latitudes ranging from Finland in the north to South Africa in the south. Limousin breeders' associations exist in many of these countries, of which 29 are members of the International Limousin Council (ILC). The ILC was founded at Limoges in 1973 by Louis de Neuville, the Limousin breed ambassador. In 1989, EUROLIM was formed to bring together all of the herd books of European Limousins. Limousins in different countries are bred according to different selection practices or objectives, and are connected by limited international gene flows. Poor genetic connectedness between countries has negative implications for estimation accuracies of international genetic prediction programmes. As a result of genetic drift or different selection, each country's population of Limousins is becoming genetically differentiated, but which is counterbalanced to a limited extent by gene flows from other countries. A study of over 2.4million Limousin pedigree files of five European countries (France, Denmark, Ireland, Sweden, United Kingdom) showed moderate gene flows from France to the United Kingdom and Denmark, but negligible gene flows to Sweden. Except for gene flows originating from France, and some limited gene flows between Denmark and Sweden in the 1990s, bull and semen exchanges between European countries has been scarce, especially since about 2000. Cow and embryo flows have been even more scarce. Conversely, the genetic contribution of North American Limousins to European countries has increased since the late 1990s, which has occurred because of their use in breeding programmes to introduce the polled gene. International Limousin genetics are now widely available in many countries for use in artificial insemination programmes, which has been facilitated by a large number of companies that specialise in the export and import of semen. Details of semen are published in extensively distributed catalogues.	Biology and health sciences	Cattle	Animals
2953473	https://en.wikipedia.org/wiki/Food%20group	Food group	Food groups categorise foods for educational purposes, usually grouping together foods with similar nutritional properties or biological classifications. Food groups are often used in nutrition guides, although the number of groups used can vary widely. Food groups were a public health education concept invented to teach people eating very restricted, unvaried diets how to avoid becoming deficient in specific nutrients. They have since been adapted to also address diseases of affluence related to diet, such as obesity, diabetes and heart disease. Historical food groups Opson and sitos were Classical Greek food groups, mainly used for moral education, to teach sophrosyne. Mitahara, a concept of moderate diet found in early-first-millennium Sanskrit texts, categorizes food into groups and recommends eating a variety of healthy foods, while avoiding the unhealthy ones; it also considers foods to have emotional and moral effects. Indian foodways had a substantial influence on European organisations such as the Vegetarian Society], which cited Indian diets as proof that a healthy vegetarian diet was possible, and were actively involved in public debate on nutrition. In the 20th century, food groups became widely used in public health education, as a tool to reduce nutritional deficiencies. As early as the 1980s, researchers were criticizing food groups, saying that they were a concept useful for teaching people to avoid nutritional deficiencies, but that nutritional deficiencies were no longer major cause of diet-related disease in affluent societies. Since these are caused by unhealthy food, not by diets lacking of a specific nutrient, they thought that food groups would have to be entirely discarded, or entirely revamped to make them useful in nutritional education in post-industrial countries. United States The USDA promoted eight basic food groups prior to 1943, then seven basic food groups until 1956, then four food groups. A food pyramid was introduced in 1992, then MyPyramid in 2005, followed by MyPlate in 2011. Dietary guidelines were introduced in 2015 and slated to be rereleased every five years. The 2020 guidelines were to be released in spring 2020. Recommended Dietary Allowance recommends daily servings of each group for a healthy diet. In the United States for instance, the USDA has described food as being in from 4 to 11 different groups. The most common food groups Dairy, also called milk products and sometimes categorized with milk alternatives or meat, is typically a smaller category in nutrition guides, if present at all, and is sometimes listed apart from other food groups. Examples of dairy products include milk, butter, ghee, yogurt, cheese, cream and ice cream. The categorization of dairy as a food group with recommended daily servings has been criticized by, for example, the Harvard School of Public Health who point out that "research has shown little benefit, and considerable potential for harm, of such high dairy intakes. Moderate consumption of milk or other dairy products—one to two servings a day—is fine, and likely has some benefits for children. But it’s not essential for adults, for a host of reasons." Fruits, sometimes categorized with vegetables, include apples, oranges, bananas, berries and lemons. Fruits contain carbohydrates, mostly in the form of non-free sugar, as well as important vitamins and minerals. Cereals and legumes, sometimes categorized as grains, is often the largest category in nutrition guides. Cereal examples include wheat, rice, oats, barley, bread and pasta. Legumes are also known as pulses and include beans, soy beans, lentils and chickpeas. Cereals are a good source of starch and are often categorized with other starchy food such as potatoes. Legumes are good source of essential amino acids as well as carbohydrates. Meat, sometimes labelled protein and occasionally inclusive of legumes and beans, eggs, meat analogues and/or dairy, is typically a medium- to smaller-sized category in nutrition guides. Examples include chicken, fish, turkey, pork and beef. Confections, also called sugary foods and sometimes categorized with fats and oils, is typically a very small category in nutrition guides, if present at all, and is sometimes listed apart from other food groups. Examples include candy, soft drinks, and chocolate. Vegetables, sometimes categorized with fruit and occasionally inclusive of legumes, is typically a large category second only to grains, or sometimes equal or superior to grains, in nutrition guides. Examples include spinach, carrots, onions, and broccoli. Water is treated in very different ways by different food guides. Some exclude the category, others list it separately from other food groups, and yet others make it the center or foundation of the guide. Water is sometimes categorized with tea, fruit juice, vegetable juice and even soup, and is typically recommended in plentiful amounts. Uncommon food groups The number of "common" food groups varies depending on who is defining them. Canada's Food Guide, which has been in continual publication since 1942 and is the second most requested government document after the income tax form in Canada, recognizes only four official food groups, listing the remainder of foods as "another".	Biology and health sciences	Health and fitness	null
2955396	https://en.wikipedia.org/wiki/Ethylenediamine	Ethylenediamine	Ethylenediamine (abbreviated as en when a ligand) is the organic compound with the formula C2H4(NH2)2. This colorless liquid with an ammonia-like odor is a basic amine. It is a widely used building block in chemical synthesis, with approximately 500,000 tonnes produced in 1998. Ethylenediamine is the first member of the so-called polyethylene amines. Synthesis Ethylenediamine is produced industrially by treating 1,2-dichloroethane with ammonia under pressure at 180 °C in an aqueous medium: In this reaction hydrogen chloride is generated, which forms a salt with the amine. The amine is liberated by addition of sodium hydroxide and can then be recovered by fractional distillation. Diethylenetriamine (DETA) and triethylenetetramine (TETA) are formed as by-products. Another industrial route to ethylenediamine involves the reaction of ethanolamine and ammonia: This process involves passing the gaseous reactants over a bed of nickel heterogeneous catalysts. It can be prepared in the lab by the reaction of either ethylene glycol or ethanolamine and urea, followed by decarboxylation of the ethyleneurea intermediate. Ethylenediamine can be purified by treatment with sodium hydroxide to remove water followed by distillation. Applications Ethylenediamine is used in large quantities for production of many industrial chemicals. It forms derivatives with carboxylic acids (including fatty acids), nitriles, alcohols (at elevated temperatures), alkylating agents, carbon disulfide, and aldehydes and ketones. Because of its bifunctional nature, having two amino groups, it readily forms heterocycles such as imidazolidines. Precursor to chelation agents, drugs, and agrochemicals A most prominent derivative of ethylenediamine is the chelating agent EDTA, which is derived from ethylenediamine via a Strecker synthesis involving cyanide and formaldehyde. Hydroxyethylethylenediamine is another commercially significant chelating agent. Numerous bio-active compounds and drugs contain the N–CH2–CH2–N linkage, including some antihistamines. Salts of ethylenebisdithiocarbamate are commercially significant fungicides under the brand names Maneb, Mancozeb, Zineb, and Metiram. Some imidazoline-containing fungicides are derived from ethylenediamine. Pharmaceutical ingredient Ethylenediamine is an ingredient in the common bronchodilator drug aminophylline, where it serves to solubilize the active ingredient theophylline. Ethylenediamine has also been used in dermatologic preparations, but has been removed from some because of causing contact dermatitis. When used as a pharmaceutical excipient, after oral administration its bioavailability is about 0.34, due to a substantial first-pass effect. Less than 20% is eliminated by renal excretion. Ethylenediamine-derived antihistamines are the oldest of the five classes of first-generation antihistamines, beginning with piperoxan aka benodain, discovered in 1933 at the Pasteur Institute in France, and also including mepyramine, tripelennamine, and antazoline. The other classes are derivatives of ethanolamine, alkylamine, piperazine, and others (primarily tricyclic and tetracyclic compounds related to phenothiazines, tricyclic antidepressants, as well as the cyproheptadine-phenindamine family) Role in polymers Ethylenediamine, because it contains two amine groups, is a widely used precursor to various polymers. Condensates derived from formaldehyde are plasticizers. It is widely used in the production of polyurethane fibers. The PAMAM class of dendrimers are derived from ethylenediamine. Tetraacetylethylenediamine The bleaching activator tetraacetylethylenediamine is generated from ethylenediamine. The derivative N,N-ethylenebis(stearamide) (EBS) is a commercially significant mold-release agent and a surfactant in gasoline and motor oil. Other applications as a solvent, it is miscible with polar solvents and is used to solubilize proteins such as albumins and casein. It is also used in certain electroplating baths. as a corrosion inhibitor in paints and coolants. ethylenediamine dihydroiodide (EDDI) is added to animal feeds as a source of iodide. chemicals for color photography developing, binders, adhesives, fabric softeners, curing agents for epoxies, and dyes. as a compound to sensitize nitromethane into an explosive. This mixture was used at Picatinny Arsenal during World War II, giving the nitromethane and ethylenediamine mixture the nickname PLX, or Picatinny Liquid Explosive. Coordination chemistry Ethylenediamine is a well-known bidentate chelating ligand for coordination compounds, with the two nitrogen atoms donating their lone pairs of electrons when ethylenediamine acts as a ligand. It is often abbreviated "en" in inorganic chemistry. The complex [Co(en)3]3+ is a well studied example. Schiff base ligands easily form from ethylenediamine. For example, the diamine condenses with 4-Trifluoromethylbenzaldehyde to give to the diimine. The salen ligands, some of which are used in catalysis, are derived from the condensation of salicylaldehydes and ethylenediamine. Related ligands Related derivatives of ethylenediamine include ethylenediaminetetraacetic acid (EDTA), tetramethylethylenediamine (TMEDA), and tetraethylethylenediamine (TEEDA). Chiral analogs of ethylenediamine include 1,2-diaminopropane and trans-diaminocyclohexane. Safety Ethylenediamine, like ammonia and other low-molecular weight amines, is a skin and respiratory irritant. Unless tightly contained, liquid ethylenediamine will release toxic and irritating vapors into its surroundings, especially on heating. The vapors absorb moisture from humid air to form a characteristic white mist, which is extremely irritating to skin, eyes, lungs and mucous membranes.	Physical sciences	Amides and amines	Chemistry
2955810	https://en.wikipedia.org/wiki/Glomeromycota	Glomeromycota	Glomeromycota (often referred to as glomeromycetes, as they include only one class, Glomeromycetes) are one of eight currently recognized divisions within the kingdom Fungi, with approximately 230 described species. Members of the Glomeromycota form arbuscular mycorrhizas (AMs) with the thalli of bryophytes and the roots of vascular land plants. Not all species have been shown to form AMs, and one, Geosiphon pyriformis, is known not to do so. Instead, it forms an endocytobiotic association with Nostoc cyanobacteria. The majority of evidence shows that the Glomeromycota are dependent on land plants (Nostoc in the case of Geosiphon) for carbon and energy, but there is recent circumstantial evidence that some species may be able to lead an independent existence. The arbuscular mycorrhizal species are terrestrial and widely distributed in soils worldwide where they form symbioses with the roots of the majority of plant species (>80%). They can also be found in wetlands, including salt-marshes, and associated with epiphytic plants. According to multigene phylogenetic analyses, this taxon is located as a member of the phylum Mucoromycota. Currently, the phylum name Glomeromycota is invalid, and the subphylum Glomeromycotina should be used to describe this taxon. Reproduction The Glomeromycota have generally coenocytic (occasionally sparsely septate) mycelia and reproduce asexually through blastic development of the hyphal tip to produce spores (Glomerospores,blastospore) with diameters of 80–500 μm. In some, complex spores form within a terminal saccule. Recently it was shown that Glomus species contain 51 genes encoding all the tools necessary for meiosis. Based on these and related findings, it was suggested that Glomus species may have a cryptic sexual cycle. Colonization New colonization of AM fungi largely depends on the amount of inoculum present in the soil. Although pre-existing hyphae and infected root fragments have been shown to colonize the roots of a host successfully, germinating spores are considered to be the key players in new host establishment. Spores are commonly dispersed by fungal and plant burrowing herbivore partners, but some air dispersal capabilities are also known. Studies have shown that spore germination is specific to particular environmental conditions such as right amount of nutrients, temperature or host availability. It has also been observed that the rate of root system colonization is directly correlated to spore density in the soil. In addition, new data also suggests that AM fungi host plants also secrete chemical factors that attract and enhance the growth of developing spore hyphae towards the root system. The necessary components for the colonization of Glomeromycota include the host's fine root system, proper development of intracellular arbuscular structures, and a well-established external fungal mycelium. Colonization is accomplished by the interactions between germinating spore hyphae and the root hairs of the host or by the development of appressoria between epidermal root cells. The process is regulated by specialized chemical signaling and changes in gene expression of both the host and AM fungi. Intracellular hyphae extend up to the cortical cells of the root and penetrate the cell walls but not the inner cellular membrane creating an internal invagination. The penetrating hyphae develop a highly branched structure called an arbuscule, which has low functional periods before degradation and absorption by the host's root cells. A fully developed arbuscular mycorrhizal structure facilitates the two-way movement of nutrients between the host and mutualistic fungal partner. The symbiotic association allows the host plant to respond better to environmental stresses, and the non-photosynthetic fungi to obtain carbohydrates produced by photosynthesis. Phylogeny Initial studies of the Glomeromycota were based on the morphology of soil-borne sporocarps (spore clusters) found in or near colonized plant roots. Distinguishing features such as wall morphologies, size, shape, color, hyphal attachment and reaction to staining compounds allowed a phylogeny to be constructed. Superficial similarities led to the initial placement of genus Glomus in the unrelated family Endogonaceae. Following broader reviews that cleared up the sporocarp confusion, the Glomeromycota were first proposed in the genera Acaulospora and Gigaspora before being accorded their own order with the three families Glomaceae (now Glomeraceae), Acaulosporaceae and Gigasporaceae. With the advent of molecular techniques this classification has undergone major revision. An analysis of small subunit (SSU) rRNA sequences indicated that they share a common ancestor with the Dikarya. Nowadays it is accepted that Glomeromycota consists of 4 orders. Several species which produce glomoid spores (i.e. spores similar to Glomus) in fact belong to other deeply divergent lineages and were placed in the orders, Paraglomerales and Archaeosporales. This new classification includes the Geosiphonaceae, which presently contains one fungus (Geosiphon pyriformis) that forms endosymbiotic associations with the cyanobacterium Nostoc punctiforme and produces spores typical to this division, in the Archaeosporales. Work in this field is incomplete, and members of Glomus may be better suited to different genera or families. Molecular biology The biochemical and genetic characterization of the Glomeromycota has been hindered by their biotrophic nature, which impedes laboratory culturing. This obstacle was eventually surpassed with the use of root cultures and, most recently, a method which applies sequencing of single nucleus from spores has also been developed to circumvent this challenge. The first mycorrhizal gene to be sequenced was the small-subunit ribosomal RNA (SSU rRNA). This gene is highly conserved and commonly used in phylogenetic studies so was isolated from spores of each taxonomic group before amplification through the polymerase chain reaction (PCR). A metatranscriptomic survey of the Sevilleta Arid Lands found that 5.4% of the fungal rRNA reads mapped to Glomeromycota. This result was inconsistent with previous PCR-based studies of community structure in the region, suggesting that previous PCR-based studies may have underestimated Glomeromycota abundance due to amplification biases.	Biology and health sciences	Basics	Plants
2956372	https://en.wikipedia.org/wiki/Foot%20per%20second	Foot per second	The foot per second (plural feet per second) is a unit of both speed (scalar) and velocity (vector quantity, which includes direction). It expresses the distance in feet (ft) traveled or displaced, divided by the time in seconds (s). The corresponding unit in the International System of Units (SI) is the meter per second. Abbreviations include ft/s, fps, and the scientific notation ft s−1. Conversions	Physical sciences	Speed	Basics and measurement
4061767	https://en.wikipedia.org/wiki/Heaviside%E2%80%93Lorentz%20units	Heaviside–Lorentz units	Heaviside–Lorentz units (or Lorentz–Heaviside units) constitute a system of units and quantities that extends the CGS with a particular set of equations that defines electromagnetic quantities, named for Oliver Heaviside and Hendrik Antoon Lorentz. They share with the CGS-Gaussian system that the electric constant and magnetic constant do not appear in the defining equations for electromagnetism, having been incorporated implicitly into the electromagnetic quantities. Heaviside–Lorentz units may be thought of as normalizing and , while at the same time revising Maxwell's equations to use the speed of light instead. The Heaviside–Lorentz unit system, like the International System of Quantities upon which the SI system is based, but unlike the CGS-Gaussian system, is rationalized, with the result that there are no factors of appearing explicitly in Maxwell's equations. That this system is rationalized partly explains its appeal in quantum field theory: the Lagrangian underlying the theory does not have any factors of when this system is used. Consequently, electromagnetic quantities in the Heaviside–Lorentz system differ by factors of in the definitions of the electric and magnetic fields and of electric charge. It is often used in relativistic calculations, and are used in particle physics. They are particularly convenient when performing calculations in spatial dimensions greater than three such as in string theory. Motivation In the mid-late 19th century, electromagnetic measurements were frequently made in either the so-named electrostatic (ESU) or electromagnetic (EMU) systems of units. These were based respectively on Coulomb's and Ampere's Law. Use of these systems, as with to the subsequently developed Gaussian CGS units, resulted in many factors of appearing in formulas for electromagnetic results, including those without any circular or spherical symmetry. For example, in the CGS-Gaussian system, the capacitance of sphere of radius is while that of a parallel plate capacitor is , where is the area of the smaller plate and is their separation. Heaviside, who was an important, though somewhat isolated, early theorist of electromagnetism, suggested in 1882 that the irrational appearance of in these sorts of relations could be removed by redefining the units for charges and fields. In his 1893 book Electromagnetic Theory, Heaviside wrote in the introduction: Length–mass–time framework As in the Gaussian system (), the Heaviside–Lorentz system () uses the length–mass–time dimensions. This means that all of the units of electric and magnetic quantities are expressible in terms of the units of the base quantities length, time and mass. Coulomb's equation, used to define charge in these systems, is in the Gaussian system, and in the HL system. The unit of charge then connects to , where 'HLC' is the HL unit of charge. The HL quantity describing a charge is then times larger than the corresponding Gaussian quantity. There are comparable relationships for the other electromagnetic quantities (see below). The commonly used set of units is the called the SI, which defines two constants, the vacuum permittivity () and the vacuum permeability (). These can be used to convert SI units to their corresponding Heaviside–Lorentz values, as detailed below. For example, SI charge is . When one puts , , , and , this evaluates to , the SI-equivalent of the Heaviside–Lorentz unit of charge. Comparison of Heaviside–Lorentz with other systems of units This section has a list of the basic formulas of electromagnetism, given in the SI, Heaviside–Lorentz, and Gaussian systems. Here and are the electric field and displacement field, respectively, and are the magnetic fields, is the polarization density, is the magnetization, is charge density, is current density, is the speed of light in vacuum, is the electric potential, is the magnetic vector potential, is the Lorentz force acting on a body of charge and velocity , is the permittivity, is the electric susceptibility, is the magnetic permeability, and is the magnetic susceptibility. Maxwell's equations The electric and magnetic fields can be written in terms of the potentials and . The definition of the magnetic field in terms of , , is the same in all systems of units, but the electric field is in the SI system, but in the HL or Gaussian systems. Other basic laws Dielectric and magnetic materials Below are the expressions for the macroscopic fields , , and in a material medium. It is assumed here for simplicity that the medium is homogeneous, linear, isotropic, and nondispersive, so that the susceptibilities are constants. Note that The quantities , and are dimensionless, and they have the same numeric value. By contrast, the electric susceptibility is dimensionless in all the systems, but has for the same material: The same statements apply for the corresponding magnetic quantities. Advantages and disadvantages of Heaviside–Lorentz units Advantages The formulas above are clearly simpler in units compared to either or Gaussian units. As Heaviside proposed, removing the from the Gauss law and putting it in the Force law considerably reduces the number of places the appears compared to Gaussian CGS units. Removing the explicit from the Gauss law makes it clear that the inverse-square force law arises by the field spreading out over the surface of a sphere. This allows a straightforward extension to other dimensions. For example, the case of long, parallel wires extending straight in the direction can be considered a two-dimensional system. Another example is in string theory, where more than three spatial dimensions often need to be considered. The equations are free of the constants and that are present in the SI system. (In addition and are overdetermined, because .) The below points are true in both Heaviside–Lorentz and Gaussian systems, but not SI. The electric and magnetic fields and have the same dimensions in the Heaviside–Lorentz system, meaning it is easy to recall where factors of go in the Maxwell equation. Every time derivative comes with a , which makes it dimensionally the same as a space derivative. In contrast, in SI units is . Giving the and fields the same dimension makes the assembly into the electromagnetic tensor more transparent. There are no factors of that need to be inserted when assembling the tensor out of the three-dimensional fields. Similarly, and have the same dimensions and are the four components of the 4-potential. The fields , , , and also have the same dimensions as and . For vacuum, any expression involving can simply be recast as the same expression with . In SI units, and have the same units, as do and , but they have different units from each other and from and . Disadvantages Despite Heaviside's urgings, it proved difficult to persuade people to switch from the established units. He believed that if the units were changed, "[o]ld style instruments would very soon be in a minority, and then disappear ...". Persuading people to switch was already difficult in 1893, and in the meanwhile there have been more than a century's worth of additional textbooks printed and voltmeters built. Heaviside–Lorentz units, like the Gaussian CGS units by which they generally differ by a factor of about 3.5, are frequently of rather inconvenient sizes. The ampere (coulomb/second) is reasonable unit for measuring currents commonly encountered, but the ESU/s, as demonstrated above, is far too small. The Gaussian CGS unit of electric potential is named a statvolt. It is about , a value which is larger than most commonly encountered potentials. The henry, the SI unit for inductance is already on the large side compared to most inductors; the Gaussian unit is 12 orders of magnitude larger. A few of the Gaussian CGS units have names; none of the Heaviside–Lorentz units do. Textbooks in theoretical physics use Heaviside–Lorentz units nearly exclusively, frequently in their natural form (see below), system's conceptual simplicity and compactness significantly clarify the discussions, and it is possible if necessary to convert the resulting answers to appropriate units after the fact by inserting appropriate factors of and . Some textbooks on classical electricity and magnetism have been written using Gaussian CGS units, but recently some of them have been rewritten to use SI units. Outside of these contexts, including for example magazine articles on electric circuits, Heaviside–Lorentz and Gaussian CGS units are rarely encountered. Translating formulas between systems To convert any formula between the SI, Heaviside–Lorentz system or Gaussian system, the corresponding expressions shown in the table below can be equated and hence substituted for each other. Replace by or vice versa. This will reproduce any of the specific formulas given in the list above. As an example, starting with the equation and the equations from the table Moving the factor across in the latter identities and substituting, the result is which then simplifies to	Physical sciences	Measurement systems	Basics and measurement
4063480	https://en.wikipedia.org/wiki/Mudskipper	Mudskipper	Mudskippers are any of the 23 extant species of amphibious fish from the subfamily Oxudercinae of the goby family Oxudercidae. They are known for their unusual body shapes, preferences for semiaquatic habitats, limited terrestrial locomotion and jumping, and the ability to survive prolonged periods of time both in and out of water. Mudskippers can grow up to long, and most are a brownish green colour that ranges anywhere from dark to light. During mating seasons, the males will also develop brightly coloured spots in order to attract females, which can be red, green or blue. Unlike other fish, the mudskipper's eyes protrude from the top of its flat head. Their most noticeable feature however is their side pectoral fins that are located more forward and under their elongated body. These fins are jointed and function similarly to limbs, which allow the mudskipper to crawl from place to place. Although having the typical body form of any other gobiid fish, these front fins allow the mudskipper to actively "skip" across muddy surfaces (hence the common name) and even climb low-hanging tree branches and scrubs. Mudskippers have also been found to be able to leap distances of up to by laterally flexing and pushing with their tails. Taxonomy Oxudercinae is sometimes classified within the family Gobiidae (gobies). Recent molecular studies do not support this classification, as oxudercine gobies appear to be paraphyletic relative to amblyopine gobies (Gobiidae: Amblyopinae), thus being included in a distinct "Periophthalmus lineage", together with amblyopines. Mudskippers can be defined as oxudercine gobies that are "fully terrestrial for some portion of the daily cycle" (character 24 in Murdy, 1989). This would define the species of the genera Boleophthalmus, Periophthalmodon, Periophthalmus, and Scartelaos as "mudskippers". However, field observations of Zappa confluentus suggest that this monotypic genus should be included in the definition. Behaviour Mudskippers typically live in burrows in intertidal habitats, and exhibit unique adaptations to this environment that are not found in most intertidal fishes, which typically survive the retreat of the tide by hiding under wet seaweed or in tide pools. These burrows are most often characterised by their smooth and vaulted ceilings. The way the males dig these burrows has been found to be directly linked to their ability to survive submerged in almost anoxic water. It has also been found to play a crucial role in the development of the eggs within the burrow. Mudskippers are quite active when out of water, feeding and interacting with one another, for example, to defend their territories and court potential partners. Once the male has completed digging his burrow he will resurface and will begin attempting to attract a female through assorted yet typical displays. These displays consist of body undulations, different postures and energetic movements. Once the female has made her choice she will then proceed to follow the male into the burrow where she will lay hundreds of eggs and allow them to be fertilized. After fertilization occurs, the period of cohabitation between the male and female is rather short. Eventually, the female will leave and it is the male that ends up guarding the egg filled burrow from predators. Mudskippers are amphibious. When leaving the water and moving into a more dry environment on land, they are still able to breathe using water that is trapped inside their large gill chambers. They are also able to absorb oxygen from the lining of their mouth and throat, allowing them to stay out of water for long periods of time. It has been discovered that they spend up to three quarters of their life on land. They are found in tropical, subtropical, and temperate regions, including the Indo-Pacific and the Atlantic coast of Africa. Adaptations Compared with fully aquatic gobies, these specialised fish present a range of anatomical and ethological adaptations that allow them to move effectively on land as well as in the water. Terrestrial movement As their name implies, these fish use their fins to move around in a series of skips. Breathing Mudskippers have the ability to breathe through their skin and the lining of their mouth (the mucosa) and throat (the pharynx); this is only possible when the mudskippers are wet, limiting them to humid habitats and requiring that they keep themselves moist. The ability to breathe through their skin is associated with increased capillary density in their skin. This mode of breathing, similar to that employed by amphibians, is known as cutaneous respiration. Another important adaptation that aids breathing while out of water is their enlarged gill chambers, where they retain a bubble of air. These chambers close tightly when the fish is above water, due to a ventromedial valve of the gill slit, keeping the gills moist, and allowing them to function while exposed to air. Gill filaments are stiff and do not coalesce when out of water. Diet The different species have adapted to various diets on the mudflats. Boleophthalmus boddarti is detritivorous, while others will eat small crabs, insects, snails and even other mudskippers. Burrowing Digging deep burrows in soft sediments allows the fish to thermoregulate, avoid marine predators during the high tide when the fish and burrow are submerged, and lay their eggs. When the burrow is submerged, several mudskipper species maintain an air pocket inside it, which allows them to breathe in conditions of very low oxygen concentration. Ammonia regulation To reduce toxic ammonia production, mudskippers can suppress amino acid breakdown when exposed to air. Another method they use involves the partial breakdown of amino acids leading to the production of alanine, which also reduces ammonia production. Mudskippers can reduce the membrane permeability of their skin and acidify the water in their burrows to reduce levels of ammonia from the environment. Blinking Mudskippers evolved the ability to blink independently from terrestrial tetrapods. Their eyes are located high on their head compared to other gobies, and they blink by lowering their eyes as a membrane called the dermal cup rises to cover them. Although other fully aquatic goby species do not have the ability to blink, mudskippers did not evolve different muscles or glands to blink with; their blinking is accomplished with the same muscles in a different configuration. Rather than having specialized glands to produce eye lubricant, the fluid film on their cornea is likely composed of mucus secreted by their skin and water from their environment, possibly stored in the infraorbital space behind the dermal cup membrane. Mudskippers likely evolved blinking in response to conditions of terrestrial life, such as to keep their eyes wet out of water (they blink more frequently in high evaporation conditions and only when colliding with things in water) and to clean and protect the eye from debris, which may adhere to the eye or approach at a faster, more dangerous speed when in air versus when in water. Their eyes are not elevated when they are still juveniles, which are fully aquatic. Species The genus Periophthalmus is by far the most diverse and widespread genus of mudskipper. Eighteen species have been described. Periophthalmus argentilineatus is one of the most widespread and well-known species. It can be found in mangrove ecosystems and mudflats of East Africa and Madagascar east through the Sundarbans of Bengal, Southeast Asia to Northern Australia, southeast China, Taiwan, and the Ryukyus, to Samoa and Tonga Islands. It grows to a length of about 9.5 cm and is a carnivorous opportunist feeder. It feeds on small prey such as small crabs and other arthropods. However, a recent molecular study suggests that P. argentilineatus is in fact a complex of species, with at least three separate lineages, one in East Africa, and two in the Indo-Malayan region. Another species, Periophthalmus barbarus, is the only oxudercine goby that inhabits the coastal areas of western Africa.	Biology and health sciences	Acanthomorpha	Animals
4064886	https://en.wikipedia.org/wiki/Hyneria	Hyneria	Hyneria is a genus of large prehistoric predatory lobe-finned fish which lived in fresh water during the Famennian stage of the Devonian period. Etymology The genus name Hyneria is a reference to the village of Hyner, Pennsylvania, near where the first specimen was found. The species epithet H. lindae is derived from the name of the wife of Keith Stewart Thomson, who described this fish. Description Hyneria was a large fish. H. lindae is estimated around in total length. An isolated cleithrum AM 6545 likely belongs to an individual of at least in length. The largest complete jaw reaches , but there is much larger fragment possibly from a jaw about twice that length, although that specimen may belong to a rhizodont instead. Assuming this jaw fragment does pertain to Hyneria, and assuming proportions similar to more complete tristichopterids, it suggests H. lindae could possibly reach lengths up to 3.5 metres (11 ft). A second species, H. udlezinye, was once estimated as having a length of between before being described. However, the species description estimates that the largest specimen belongs to an animal about . Its skull had heavy, ornamented dermal bones and its lower jaw was relatively long and shallow. The teeth were stout, with those of the premaxilla forming fangs upwards of . Its body was covered by cycloid scales. It had large sensory canals to aid in detection of possible prey, as the freshwater environment it inhabited likely was murky and had low visibility. Adult individuals retained juvenile features (i.e. partially unossified skeletons), suggesting that they were likely neotenic. Discovery The original fossils came from two localities in Pennsylvania, United States, one found between the villages of North Bend and Hyner and another near Emporium. They consisted of a disarticulated partial skull and fragments of the shoulder girdle. The fossils were found in the Catskill Formation of the Red Hill Shale, dating to the upper Devonian. These were the only remains known until 1993 when a renewed collecting effort discovered abundant new material. Hyneria is considered the largest and most common lobe-finned fish found in the Red Hill Shale. In February 2023 a second species of Hyneria, H. udlezinye, was named from remains discovered in the Waterloo Farm lagerstätte. These remains include the skull and shoulder girdle.	Biology and health sciences	Prehistoric osteichthyans	Animals
4068447	https://en.wikipedia.org/wiki/Successive%20over-relaxation	Successive over-relaxation	In numerical linear algebra, the method of successive over-relaxation (SOR) is a variant of the Gauss–Seidel method for solving a linear system of equations, resulting in faster convergence. A similar method can be used for any slowly converging iterative process. It was devised simultaneously by David M. Young Jr. and by Stanley P. Frankel in 1950 for the purpose of automatically solving linear systems on digital computers. Over-relaxation methods had been used before the work of Young and Frankel. An example is the method of Lewis Fry Richardson, and the methods developed by R. V. Southwell. However, these methods were designed for computation by human calculators, requiring some expertise to ensure convergence to the solution which made them inapplicable for programming on digital computers. These aspects are discussed in the thesis of David M. Young Jr. Formulation Given a square system of n linear equations with unknown x: where: Then A can be decomposed into a diagonal component D, and strictly lower and upper triangular components L and U: where The system of linear equations may be rewritten as: for a constant ω > 1, called the relaxation factor. The method of successive over-relaxation is an iterative technique that solves the left hand side of this expression for x, using the previous value for x on the right hand side. Analytically, this may be written as: where is the kth approximation or iteration of and is the next or k + 1 iteration of . However, by taking advantage of the triangular form of (D+ωL), the elements of x(k+1) can be computed sequentially using forward substitution: This can again be written analytically in matrix-vector form without the need of inverting the matrix : Convergence The choice of relaxation factor ω is not necessarily easy, and depends upon the properties of the coefficient matrix. In 1947, Ostrowski proved that if is symmetric and positive-definite then for . Thus, convergence of the iteration process follows, but we are generally interested in faster convergence rather than just convergence. Convergence Rate The convergence rate for the SOR method can be analytically derived. One needs to assume the following the relaxation parameter is appropriate: Jacobi's iteration matrix has only real eigenvalues Jacobi's method is convergent: the matrix decomposition satisfies the property that for any and . Then the convergence rate can be expressed as where the optimal relaxation parameter is given by In particular, for (Gauss-Seidel) it holds that . For the optimal we get , which shows SOR is roughly four times more efficient than Gauss–Seidel. The last assumption is satisfied for tridiagonal matrices since for diagonal with entries and . Algorithm Since elements can be overwritten as they are computed in this algorithm, only one storage vector is needed, and vector indexing is omitted. The algorithm goes as follows: Inputs: , , Output: Choose an initial guess to the solution repeat until convergence for from 1 until do set to 0 for from 1 until do if ≠ then set to end if end (-loop) set to end (-loop) check if convergence is reached end (repeat) Note can also be written , thus saving one multiplication in each iteration of the outer for-loop. Example We are presented the linear system To solve the equations, we choose a relaxation factor and an initial guess vector . According to the successive over-relaxation algorithm, the following table is obtained, representing an exemplary iteration with approximations, which ideally, but not necessarily, finds the exact solution, , in 38 steps. A simple implementation of the algorithm in Common Lisp is offered below. ;; Set the default floating-point format to "long-float" in order to ;; ensure correct operation on a wider range of numbers. (setf read-default-float-format 'long-float) (defparameter +MAXIMUM-NUMBER-OF-ITERATIONS+ 100 "The number of iterations beyond which the algorithm should cease its operation, regardless of its current solution. A higher number of iterations might provide a more accurate result, but imposes higher performance requirements.") (declaim (type (integer 0 ) +MAXIMUM-NUMBER-OF-ITERATIONS+)) (defun get-errors (computed-solution exact-solution) "For each component of the COMPUTED-SOLUTION vector, retrieves its error with respect to the expected EXACT-SOLUTION vector, returning a vector of error values. --- While both input vectors should be equal in size, this condition is not checked and the shortest of the twain determines the output vector's number of elements. --- The established formula is the following: Let resultVectorSize = min(computedSolution.length, exactSolution.length) Let resultVector = new vector of resultVectorSize For i from 0 to (resultVectorSize - 1) resultVector[i] = exactSolution[i] - computedSolution[i] Return resultVector" (declare (type (vector number ) computed-solution)) (declare (type (vector number ) exact-solution)) (map '(vector number ) #'- exact-solution computed-solution)) (defun is-convergent (errors &key (error-tolerance 0.001)) "Checks whether the convergence is reached with respect to the ERRORS vector which registers the discrepancy betwixt the computed and the exact solution vector. --- The convergence is fulfilled if and only if each absolute error component is less than or equal to the ERROR-TOLERANCE, that is: For all e in ERRORS, it holds: abs(e) <= errorTolerance." (declare (type (vector number ) errors)) (declare (type number error-tolerance)) (flet ((error-is-acceptable (error) (declare (type number error)) (<= (abs error) error-tolerance))) (every #'error-is-acceptable errors))) (defun make-zero-vector (size) "Creates and returns a vector of the SIZE with all elements set to 0." (declare (type (integer 0 ) size)) (make-array size :initial-element 0.0 :element-type 'number)) (defun successive-over-relaxation (A b omega &key (phi (make-zero-vector (length b))) (convergence-check #'(lambda (iteration phi) (declare (ignore phi)) (>= iteration +MAXIMUM-NUMBER-OF-ITERATIONS+)))) "Implements the successive over-relaxation (SOR) method, applied upon the linear equations defined by the matrix A and the right-hand side vector B, employing the relaxation factor OMEGA, returning the calculated solution vector. --- The first algorithm step, the choice of an initial guess PHI, is represented by the optional keyword parameter PHI, which defaults to a zero-vector of the same structure as B. If supplied, this vector will be destructively modified. In any case, the PHI vector constitutes the function's result value. --- The terminating condition is implemented by the CONVERGENCE-CHECK, an optional predicate lambda(iteration phi) => generalized-boolean which returns T, signifying the immediate termination, upon achieving convergence, or NIL, signaling continuant operation, otherwise. In its default configuration, the CONVERGENCE-CHECK simply abides the iteration's ascension to the ``+MAXIMUM-NUMBER-OF-ITERATIONS+'', ignoring the achieved accuracy of the vector PHI." (declare (type (array number (* )) A)) (declare (type (vector number ) b)) (declare (type number omega)) (declare (type (vector number ) phi)) (declare (type (function ((integer 1 ) (vector number )) ) convergence-check)) (let ((n (array-dimension A 0))) (declare (type (integer 0 ) n)) (loop for iteration from 1 by 1 do (loop for i from 0 below n by 1 do (let ((rho 0)) (declare (type number rho)) (loop for j from 0 below n by 1 do (when (/= j i) (let ((a[ij] (aref A i j)) (phi[j] (aref phi j))) (incf rho ( a[ij] phi[j]))))) (setf (aref phi i) (+ (* (- 1 omega) (aref phi i)) (* (/ omega (aref A i i)) (- (aref b i) rho)))))) (format T "~&~d. solution = ~a" iteration phi) ;; Check if convergence is reached. (when (funcall convergence-check iteration phi) (return)))) (the (vector number ) phi)) ;; Summon the function with the exemplary parameters. (let ((A (make-array (list 4 4) :initial-contents '(( 4 -1 -6 0 ) ( -5 -4 10 8 ) ( 0 9 4 -2 ) ( 1 0 -7 5 )))) (b (vector 2 21 -12 -6)) (omega 0.5) (exact-solution (vector 3 -2 2 1))) (successive-over-relaxation A b omega :convergence-check #'(lambda (iteration phi) (declare (type (integer 0 ) iteration)) (declare (type (vector number ) phi)) (let ((errors (get-errors phi exact-solution))) (declare (type (vector number ) errors)) (format T "~&~d. errors = ~a" iteration errors) (or (is-convergent errors :error-tolerance 0.0) (>= iteration +MAXIMUM-NUMBER-OF-ITERATIONS+)))))) A simple Python implementation of the pseudo-code provided above. import numpy as np from scipy import linalg def sor_solver(A, b, omega, initial_guess, convergence_criteria): """ This is an implementation of the pseudo-code provided in the Wikipedia article. Arguments: A: nxn numpy matrix. b: n dimensional numpy vector. omega: relaxation factor. initial_guess: An initial solution guess for the solver to start with. convergence_criteria: The maximum discrepancy acceptable to regard the current solution as fitting. Returns: phi: solution vector of dimension n. """ step = 0 phi = initial_guess[:] residual = linalg.norm(A @ phi - b) # Initial residual while residual > convergence_criteria: for i in range(A.shape[0]): sigma = 0 for j in range(A.shape[1]): if j != i: sigma += A[i, j] * phi[j] phi[i] = (1 - omega) * phi[i] + (omega / A[i, i]) * (b[i] - sigma) residual = linalg.norm(A @ phi - b) step += 1 print("Step {} Residual: {:10.6g}".format(step, residual)) return phi # An example case that mirrors the one in the Wikipedia article residual_convergence = 1e-8 omega = 0.5 # Relaxation factor A = np.array([[4, -1, -6, 0], [-5, -4, 10, 8], [0, 9, 4, -2], [1, 0, -7, 5]]) b = np.array([2, 21, -12, -6]) initial_guess = np.zeros(4) phi = sor_solver(A, b, omega, initial_guess, residual_convergence) print(phi) Symmetric successive over-relaxation The version for symmetric matrices A, in which is referred to as Symmetric Successive Over-Relaxation, or (SSOR), in which and the iterative method is The SOR and SSOR methods are credited to David M. Young Jr. Other applications of the method A similar technique can be used for any iterative method. If the original iteration had the form then the modified version would use However, the formulation presented above, used for solving systems of linear equations, is not a special case of this formulation if is considered to be the complete vector. If this formulation is used instead, the equation for calculating the next vector will look like where . Values of are used to speed up convergence of a slow-converging process, while values of are often used to help establish convergence of a diverging iterative process or speed up the convergence of an overshooting process. There are various methods that adaptively set the relaxation parameter based on the observed behavior of the converging process. Usually they help to reach a super-linear convergence for some problems but fail for the others.	Mathematics	Linear algebra	null
5417918	https://en.wikipedia.org/wiki/Basement%20%28geology%29	Basement (geology)	In geology, basement and crystalline basement are crystalline rocks lying above the mantle and beneath all other rocks and sediments. They are sometimes exposed at the surface, but often they are buried under miles of rock and sediment. The basement rocks lie below a sedimentary platform or cover, or more generally any rock below sedimentary rocks or sedimentary basins that are metamorphic or igneous in origin. In the same way, the sediments or sedimentary rocks on top of the basement can be called a "cover" or "sedimentary cover". Crustal rocks are modified several times before they become basement, and these transitions alter their composition. Continental crust Basement rock is the thick foundation of ancient, and oldest, metamorphic and igneous rock that forms the crust of continents, often in the form of granite. Basement rock is contrasted to overlying sedimentary rocks which are laid down on top of the basement rocks after the continent was formed, such as sandstone and limestone. The sedimentary rocks which may be deposited on top of the basement usually form a relatively thin veneer, but can be more than thick. The basement rock of the crust can be thick or more. The basement rock can be located under layers of sedimentary rock, or be visible at the surface. Basement rock is visible, for example, at the bottom of the Grand Canyon, consisting of 1.7- to 2-billion-year-old granite (Zoroaster Granite) and schist (Vishnu Schist). The Vishnu Schist is believed to be highly metamorphosed igneous rocks and shale, from basalt, mud and clay laid from volcanic eruptions, and the granite is the result of magma intrusions into the Vishnu Schist. An extensive cross section of sedimentary rocks laid down on top of it through the ages is visible as well. Age The basement rocks of the continental crust tend to be much older than the oceanic crust. The oceanic crust can be from 0–340 million years in age, with an average age of 64 million years. Continental crust is older because continental crust is light and thick enough so it is not subducted, while oceanic crust is periodically subducted and replaced at subduction and oceanic rifting areas. Complexity The basement rocks are often highly metamorphosed and complex, and are usually crystalline. They may consist of many different types of rock – volcanic, intrusive igneous and metamorphic. They may also contain ophiolites, which are fragments of oceanic crust that became wedged between plates when a terrane was accreted to the edge of the continent. Any of this material may be folded, refolded and metamorphosed. New igneous rock may freshly intrude into the crust from underneath, or may form underplating, where the new igneous rock forms a layer on the underside of the crust. The majority of continental crust on the planet is around 1 to 3 billion years old, and it is theorised that there was at least one period of rapid expansion and accretion to the continents during the Precambrian. Much of the basement rock may have originally been oceanic crust, but it was highly metamorphosed and converted into continental crust. It is possible for oceanic crust to be subducted down into the Earth's mantle, at subduction fronts, where oceanic crust is being pushed down into the mantle by an overriding plate of oceanic or continental crust. Volcanism When a plate of oceanic crust is subducted beneath an overriding plate of oceanic crust, as the underthrusting crust melts, it causes an upwelling of magma that can cause volcanism along the subduction front on the overriding plate. This produces an oceanic volcanic arc, like Japan. This volcanism causes metamorphism, introduces igneous intrusions, and thickens the crust by depositing additional layers of extrusive igneous rock from volcanoes. This tends to make the crust thicker and less dense, making it immune to subduction. Oceanic crust can be subducted, while continental crust cannot. Eventually, the subduction of the underthrusting oceanic crust can bring the volcanic arc close to a continent, with which it may collide. When this happens, instead of being subducted, it is accreted to the edge of the continent and becomes part of it. Thin strips or fragments of the underthrusting oceanic plate may also remain attached to the edge of the continent so that they are wedged and tilted between the converging plates, creating ophiolites. In this manner, continents can grow over time as new terranes are accreted to their edges, and so continents can be composed of a complex quilt of terranes of varying ages. As such, the basement rock can become younger going closer to the edge of the continent. There are exceptions, however, such as exotic terranes. Exotic terranes are pieces of other continents that have broken off from their original parent continent and have become accreted to a different continent. Cratons Continents can consist of several continental cratons – blocks of crust built around an initial original core of continents – that gradually grew and expanded as additional newly created terranes were added to their edges. For instance, Pangea consisted of most of the Earth's continents being accreted into one giant supercontinent. Most continents, such as Asia, Africa and Europe, include several continental cratons, as they were formed by the accretion of many smaller continents. Usage In European geology, the basement generally refers to rocks older than the Variscan orogeny. On top of this older basement Permian evaporites and Mesozoic limestones were deposited. The evaporites formed a weak zone on which the harder (stronger) limestone cover was able to move over the hard basement, making the distinction between basement and cover even more pronounced. In Andean geology the basement refers to the Proterozoic, Paleozoic and early Mesozoic (Triassic to Jurassic) rock units as the basement to the late Mesozoic and Cenozoic Andean sequences developed following the onset of subduction along the western margin of the South American Plate. When discussing the Trans-Mexican Volcanic Belt of Mexico the basement include Proterozoic, Paleozoic and Mesozoic age rocks for the Oaxaquia, the Mixteco and the Guerrero terranes respectively. The term basement is used mostly in disciplines of geology like basin geology, sedimentology and petroleum geology in which the (typically Precambrian) crystalline basement is not of interest as it rarely contains petroleum or natural gas. The term economic basement is also used to describe the deeper parts of a cover sequence that are of no economic interest.	Physical sciences	Stratigraphy	Earth science
963313	https://en.wikipedia.org/wiki/Hemp	Hemp	Hemp, or industrial hemp, is a plant in the botanical class of Cannabis sativa cultivars grown specifically for industrial and consumable use. It can be used to make a wide range of products. Along with bamboo, hemp is among the fastest growing plants on Earth. It was also one of the first plants to be spun into usable fiber 50,000 years ago. It can be refined into a variety of commercial items, including paper, rope, textiles, clothing, biodegradable plastics, paint, insulation, biofuel, food, and animal feed. Although chemotype I cannabis and hemp (types II, III, IV, V) are both Cannabis sativa and contain the psychoactive component tetrahydrocannabinol (THC), they represent distinct cultivar groups, typically with unique phytochemical compositions and uses. Hemp typically has lower concentrations of total THC and may have higher concentrations of cannabidiol (CBD), which potentially mitigates the psychoactive effects of THC. The legality of hemp varies widely among countries. Some governments regulate the concentration of THC and permit only hemp that is bred with an especially low THC content into commercial production. Etymology The etymology is uncertain but there appears to be no common Proto-Indo-European source for the various forms of the word; the Greek term () is the oldest attested form, which may have been borrowed from an earlier Scythian or Thracian word. Then it appears to have been borrowed into Latin, and separately into Slavic and from there into Baltic, Finnish, and Germanic languages. In the Germanic languages, following Grimm's law, the "k" would have changed to "h" with the first Germanic sound shift, giving Proto-Germanic *hanapiz, after which it may have been adapted into the Old English form, , . Barber (1991) however, argued that the spread of the name "kannabis" was due to its historically more recent plant use, starting from the south, around Iran, whereas non-THC varieties of hemp are older and prehistoric. Another possible source of origin is Assyrian , which was the name for a source of oil, fiber, and medicine in the 1st millennium BC. Cognates of hemp in other Germanic languages include Dutch , Danish and Norwegian , Saterland Frisian , German , Icelandic and Swedish . In those languages "hemp" can refer to either industrial fiber hemp or narcotic cannabis strains. Uses Hemp is used to make a variety of commercial and industrial products, including rope, textiles, clothing, shoes, food, paper, bioplastics, insulation, and biofuel. The bast fibers can be used to make textiles that are 100% hemp, but they are commonly blended with other fibers, such as flax, cotton or silk, as well as virgin and recycled polyester, to make woven fabrics for apparel and furnishings. The inner two fibers of the plant are woodier and typically have industrial applications, such as mulch, animal bedding, and litter. When oxidized (often erroneously referred to as "drying"), hemp oil from the seeds becomes solid and can be used in the manufacture of oil-based paints, in creams as a moisturizing agent, for cooking, and in plastics. Hemp seeds have been used in bird feed mix as well. A survey in 2003 showed that more than 95% of hemp seed sold in the European Union was used in animal and bird feed. Food Hemp seeds can be eaten raw, ground into hemp meal, sprouted or made into dried sprout powder. Hemp seeds can also be made into a slurry used for baking or for beverages, such as hemp milk and tisanes. Hemp oil is cold-pressed from the seed and is high in unsaturated fatty acids. In the UK, the Department for Environment, Food and Rural Affairs treats hemp as a purely non-food crop, but with proper licensing and proof of less than 0.3% THC concentration, hemp seeds can be imported for sowing or for sale as a food or food ingredient. In the US, hemp can be used legally in food products and, , was typically sold in health food stores or through mail order. Nutrition A portion of hulled hemp seeds supplies of food energy. They contain 5% water, 5% carbohydrates, 49% total fat, and 31% protein. The share of protein obtained from the hemp seeds can be increased in by processing the seeds, such as by dehulling the seeds, or by using the meal or cake (also called hemp seed flour), that is, the remaining fraction of hemp seed obtained after expelling its oil fraction. The proteins are mostly located in the inner layer of the seed, whereas the hull is poor in proteins, as it mostly contains the fiber. Hemp seeds are notable in providing 64% of the Daily Value (DV) of protein per 100-gram serving. The three main proteins in hemp seeds are edestin (83% of total protein content), albumin (13%) and ß-conglycinin (up to 5%). Hemp seed proteins are highly digestible compared to soy proteins when untreated (unheated). The amino acid profile of hemp seeds is comparable to the profiles of other protein-rich foods, such as meat, milk, eggs, and soy. Protein digestibility-corrected amino acid scores were 0.49–0.53 for whole hemp seed, 0.46–0.51 for hemp seed meal, and 0.63–0.66 for hulled hemp seed. The most abundant amino acid in hemp seed is glutamic acid (3.74–4.58% of whole seed) followed by arginine (2.28–3.10% of whole seed). The whole hemp seed can be considered a rich-protein source containing a protein amount higher or similar than other protein-rich products, such as quinoa (13.0%), chia seeds (18.2–19.7%), buckwheat seeds (27.8%) and linseeds (20.9%). Nutritionally, the protein fraction of hemp seed is highly digestible comparing to other plant-based proteins such as soy protein. Hemp seed protein has a good profile of essential amino acids, still, this profile of amino acids is inferior to that of soy or casein. Hemp seeds are a rich source of dietary fiber (20% DV), B vitamins, and the dietary minerals manganese (362% DV), phosphorus (236% DV), magnesium (197% DV), zinc (104% DV), and iron (61% DV). About 73% of the energy in hemp seeds is in the form of fats and essential fatty acids, mainly polyunsaturated fatty acids, linoleic, oleic, and alpha-linolenic acids. The ratio of the 38.100 grams of polyunsaturated fats per 100 grams is 9.301 grams of omega-3 to 28.698 grams of omega-6. Typically, the portion suggested on packages for an adult is 30 grams, approximately three tablespoons. With its gluten content as low as 4.78 ppm, hemp is attracting attention as a gluten-free (<20 ppm) food material. Despite the rich nutrient content of hemp seeds, the seeds contain antinutritional compounds, including phytic acid, trypsin inhibitors, and tannins, in statistically significant concentrations. Storage Hemp oil oxidizes and turns rancid within a short period of time if not stored properly; its shelf life is extended when it is stored in a dark airtight container and refrigerated. Both light and heat can degrade hemp oil. Fiber Hemp fiber has been used extensively throughout history, with production climaxing soon after being introduced to the New World. For centuries, items ranging from rope, to fabrics, to industrial materials were made from hemp fiber. Hemp was also commonly used to make sail canvas. The word "canvas" is derived from the word cannabis. Pure hemp has a texture similar to linen. Because of its versatility for use in a variety of products, today hemp is used in a number of consumer goods, including clothing, shoes, accessories, dog collars, and home wares. For clothing, in some instances, hemp is mixed with lyocell. Its benefits in terms for sustainability also increase its appeal in industries, such as the clothing industry. Building material Hemp as a building construction material provides solutions to a variety of issues facing current building standards. Its light weight, mold resistance, breathability, etc. makes hemp products versatile in a multitude of uses. Following the co-heating tests of NNFCC Renewable House at the Building Research Establishment (BRE), hemp is reported to be a more sustainable material of construction in comparison to most building methods used today. In addition, its practical use in building construction could result in the reduction of both energy consumption costs and the creation of secondary pollutants. In 2022, hemp-lime, also known as hempcrete, was accepted as a building material, along with methodologies for its use, by the International Code Council, and was included in the 2024 edition of the International Residential Code as an appendix: "Appendix BL Hemp-Lime (Hempcrete) Construction". This inclusion in the IRC model code is expected to promote expansion of the use and legitimacy of hemp-lime in construction in the United States. The hemp market was at its largest during the 17th century. In the 19th century and onward, the market saw a decline during its rapid illegalization in many countries. Hemp has resurfaced in green building construction, primarily in Europe. The modern-day disputes regarding the legality of hemp lead to its main disadvantages: importing and regulating costs. Final Report on the Construction of the Hemp Houses at Haverhill, UK conducts that hemp construction exceeds the cost of traditional building materials by £48per square meter. Currently, the University of Bath researches the use of hemp-lime panel systems for construction. Funded by the European Union, the research tests panel design within their use in high-quality construction, on site assembly, humidity and moisture penetration, temperature change, daily performance and energy saving documentations. The program, focusing on Britain, France, and Spain markets aims to perfect protocols of use and application, manufacturing, data gathering, certification for market use, as well as warranty and insurance. The most common use of hemp-lime in building is by casting the hemp-hurd and lime mix while wet around a timber frame with temporary shuttering and tamping the mix to form a firm mass. After the removal of the temporary shuttering, the solidified hemp mix is then ready to be plastered with lime plaster. Sustainability Hemp is classified under the green category of building design, primarily due to its positive effects on the environment. A few of its benefits include but are not limited to the suppression of weed growth, anti-erosion, reclamation properties, and the ability to remove poisonous substances and heavy metals from soil. The use of hemp is beginning to gain popularity alongside other natural materials. This is because cannabis processing is done mechanically with minimal harmful effects on the environment. A part of what makes hemp sustainable is its minimal water usage and non-reliance on pesticides for proper growth. It is recyclable, non-toxic, and biodegradable, making hemp a popular choice in green building construction. Hemp fiber is known to have high strength and durability, and has been known to be a good protector against vermin. The fiber has the capability to reinforce structures by embossing threads and cannabis shavers. Hemp has been involved more recently in the building industry, producing building construction materials including insulation, hempcrete, and varnishes. Hemp made materials have low embodied energy. The plant has the ability to absorb large amounts of CO2, providing air quality, thermal balance, creating a positive environmental impact. Hemp's properties allow mold resistance, and its porous materiality makes the building materials made of it breathable. In addition hemp possesses the ability to absorb and release moisture without deteriorating. Hemp can be non-flammable if mixed with lime and could be applied on numerous aspects of the building (wall, roofs, etc.) due to its lightweight properties. Insulation Hemp is commonly used as an insulation material. Its flexibility and toughness during compression allows for easier implementation within structural framing systems. The insulation material could also be easily adjusted to different sizes and shapes by being cut during the installation process. The ability to not settle and therefore avoiding cavity developments lowers its need for maintenance. Hemp insulation is naturally lightweight and non-toxic, allowing for an exposed installation in a variety of spaces, including flooring, walling, and roofing. Compared to mineral insulation, hemp absorbs roughly double the amount of heat and could be compared to wood, in some cases even overpassing some of its types. Hemp insulation's porous materiality allows for air and moisture penetration, with a bulk density going up to 20% without losing any thermal properties. In contrast, the commonly used mineral insulation starts to fail after 2%. The insulation evenly distributes vapor and allows for air circulation, constantly carrying out used air and replacing with fresh. Its use on the exterior of the structure, overlaid with breathable water-resistive barriers, eases the withdrawal of moisture from within the wall structure. In addition, the insulation doubles as a sound barrier, weakening airborne sound waves passing through it. Hempcrete In addition to the CO2 absorbed during its growth period, hemp-lime, also known as hempcrete, continues absorption during the curing process. The mixture hardens when the silica contained in hemp shives mixes with hydraulic lime, resulting in the mineralization process called "carbonation".. Though not a load-bearing material, hempcrete is most commonly used as infill in building construction due to its light weight (roughly seven times lighter than common concrete) and vapor permeability. The building material is made of hemp hurds (shiv or shives), hydraulic lime, and water mixed in varying ratios. The mix depends on the use of the material within the structure and could differ in physical properties. Surfaces such as flooring interact with a multitude of loads and would have to be more resistive, while walls and roofs are required to be more lightweight. The application of this material in construction requires minimal skill. Hempcrete can be formed in-situ or formed into blocks. Such blocks are not strong enough to be used for structural elements and must be supported by brick, wood, or steel framing. In the end of the twentieth century, during his renovation of Maison de la Turquie in Nogent-sur-Seine, France, Charles Rasetti first invented and applied the use of hempcrete in construction. Shortly after, in the 2000s, Modece Architects used hemp-lime for test designs in Haverhill. The dwellings were studied and monitored for comparison with other building performances by BRE. Completed in 2009, the Center for the Built Environment's Renewable House was found to be among the most technologically advanced structures made of hemp-based material. A year later the first home made of hemp-based materials was completed in Asheville, North Carolina, US. Oils and varnishes Cannabis seeds have high-fat content and contain 30-35% of fatty acids. The extracted oil is suited for a variety of construction applications. The biodegradable hemp oil acts as a wood varnish, protecting flooring from mold, pests, and wear. Its use prevents the water from penetrating the wood while still allowing air and vapor to pass through. Its most common use can be seen in wood framing construction, one of the most common construction methods in the world. Because of its low UV-resistant rating, the finish is most often used indoors, on surfaces such as flooring and wood paneling. Plaster Hemp-based insulating plaster is created by combining hemp fibers with calcium lime and sand. This material, when applied on internal walls, ceilings, and flooring, can be layered up to ten centimeters in thickness. Its porous materiality allows the created plaster to regulate air humidity and evenly distribute it. The gradual absorption and release of water prevent the material from cracking and breaking apart. Similar to high-density fiber cement, hemp plaster can naturally vary in color and be manually pigmented. Ropes and strands Hemp ropes can be woven in various diameters, possessing high amounts of strength making them suitable for a variety of uses for building construction purposes. Some of these uses include installation of frames in building openings and connection of joints. The ropes also used in bridge construction, tunnels, traditional homes, etc. One of the earliest examples of hemp rope and other textile use can be traced back to 1500 BC Egypt. Plastics Cannabis geotextiles could be put in both wet and dry conditions. Hemp-based bioplastic is a biodegradable alternative to regular plastic and can potentially replace polyvinyl chloride (PVC), a material used for plumbing pipes. Wood Hemp growth lasts roughly 100 days, a much faster time period than an average tree used for construction purposes. While dry, the fibers could be pressed into tight wood alternatives to wood-frame construction, wall/ceiling paneling, and flooring. As an addition, hemp is flexible and versatile allowing it to be used in a greater number of ways than wood. Similarly, hemp wood could also be made of recycled hemp-based paper. Composite materials A mixture of fiberglass, hemp fiber, kenaf, and flax has been used since 2002 to make composite panels for automobiles. The choice of which bast fiber to use is primarily based on cost and availability. Various car makers are beginning to use hemp in their cars, including Audi, BMW, Ford, GM, Chrysler, Honda, Iveco, Lotus, Mercedes, Mitsubishi, Porsche, Saturn, Volkswagen and Volvo. For example, the Lotus Eco Elise and the Mercedes C-Class both contain hemp (up to 20 kg in each car in the case of the latter). Paper Hemp paper are paper varieties consisting exclusively or to a large extent from pulp obtained from fibers of industrial hemp. The products are mainly specialty papers such as cigarette paper, banknotes and technical filter papers. Compared to wood pulp, hemp pulp offers a four to five times longer fiber, a significantly lower lignin fraction as well as a higher tear resistance and tensile strength. However, production costs are about four times higher than for paper from wood, since the infrastructure for using hemp is underdeveloped. If the paper industry were to switch from wood to hemp for sourcing its cellulose fibers, the following benefits could be utilized: Hemp yields three to four times more usable fiber per hectare per annum than forests, and hemp does not need pesticides or herbicides. Hemp has a much faster crop yield. It takes about 3–4 months for hemp stalks to reach maturity, while trees can take between 20 and 80 years. Not only does hemp grow at a faster rate, but it also contains a high level of cellulose. This quick return means that paper can be produced at a faster rate if hemp were used in place of wood. Hemp paper does not require the use of toxic bleaching or as many chemicals as wood pulp because it can be whitened with hydrogen peroxide. This means using hemp instead of wood for paper would end the practice of poisoning Earth's waterways with chlorine or dioxins from wood paper manufacturing. Hemp paper can be recycled up to 8 times, compared to just 3 times for paper made from wood pulp. Compared to its wood pulp counterpart, paper from hemp fibers resists decomposition and does not yellow or brown with age. It is also one of the strongest natural fibers in the world - one of the reasons for its longevity and durability. Several factors favor the increased use of wood substitutes for paper, especially agricultural fibers such as hemp. Deforestation, particularly the destruction of old growth forests, and the world's decreasing supply of wild timber resources are today major ecological concerns. Hemp's use as a wood substitute will contribute to preserving biodiversity. However, hemp has had a hard time competing with paper from trees or recycled newsprint. Only the outer part of the stem consists mainly of fibers which are suitable for the production of paper. Numerous attempts have been made to develop machines that efficiently and inexpensively separate useful fibers from less useful fibers, but none have been completely successful. This has meant that paper from hemp is still expensive compared to paper from trees. Jewelry Hemp jewelry is the product of knotting hemp twine through the practice of macramé. Hemp jewelry includes bracelets, necklaces, anklets, rings, watches, and other adornments. Some jewelry features beads made from crystals, glass, stone, wood and bones. The hemp twine varies in thickness and comes in a variety of colors. There are many different stitches used to create hemp jewelry, however, the half knot and full knot stitches are most common. Cordage Hemp rope was used in the age of sailing ships, though the rope had to be protected by tarring, since hemp rope has a propensity for breaking from rot, as the capillary effect of the rope-woven fibers tended to hold liquid at the interior, while seeming dry from the outside. Tarring was a labor-intensive process, and earned sailors the nickname "Jack Tar". Hemp rope was phased out when manila rope, which does not require tarring, became widely available. Manila is sometimes referred to as Manila hemp, but is not related to hemp; it is abacá, a species of banana. Animal bedding Hemp shives are the core of the stem, hemp hurds are broken parts of the core. In the EU, they are used for animal bedding (horses, for instance), or for horticultural mulch. Industrial hemp is much more profitable if both fibers and shives (or even seeds) can be used. Water and soil purification Hemp can be used as a "mop crop" to clear impurities out of wastewater, such as sewage effluent, excessive phosphorus from chicken litter, or other unwanted substances or chemicals. Additionally, hemp is being used to clean contaminants at the Chernobyl nuclear disaster site, by way of a process which is known as phytoremediation – the process of clearing radioisotopes and a variety of other toxins from the soil, water, and air. Weed control Hemp crops are tall, have thick foliage, and can be planted densely, and thus can be grown as a smother crop to kill tough weeds. Using hemp this way can help farmers avoid the use of herbicides, gain organic certification, and gain the benefits of crop rotation. However, due to the plant's rapid and dense growth characteristics, some jurisdictions consider hemp a prohibited and noxious weed, much like Scotch Broom. Biofuels Biodiesel can be made from the oils in hemp seeds and stalks; this product is sometimes called "hempoline". Alcohol fuel (ethanol or, less commonly, methanol) can be made by fermenting the whole plant. Filtered hemp oil can be used directly to power diesel engines. In 1892, Rudolf Diesel invented the diesel engine, which he intended to power "by a variety of fuels, especially vegetable and seed oils, which earlier were used for oil lamps, i.e. the Argand lamp". Production of vehicle fuel from hemp is very small. Commercial biodiesel and biogas is typically produced from cereals, coconuts, palm seeds, and cheaper raw materials like garbage, wastewater, dead plant and animal material, animal feces and kitchen waste. Processing Separation of hurd and bast fiber is known as decortication. Traditionally, hemp stalks would be water-retted first before the fibers were beaten off the inner hurd by hand, a process known as scutching. As mechanical technology evolved, separating the fiber from the core was accomplished by crushing rollers and brush rollers, or by hammer-milling, wherein a mechanical hammer mechanism beats the hemp against a screen until hurd, smaller bast fibers, and dust fall through the screen. After the Marijuana Tax Act was implemented in 1938, the technology for separating the fibers from the core remained "frozen in time". Recently, new high-speed kinematic decortication has come about, capable of separating hemp into three streams; bast fiber, hurd, and green microfiber. Only in 1997, did Ireland, parts of the Commonwealth and other countries begin to legally grow industrial hemp again. Iterations of the 1930s decorticator have been met with limited success, along with steam explosion and chemical processing known as thermomechanical pulping. Cultivation Hemp is usually planted between March and May in the northern hemisphere, between September and November in the southern hemisphere. It matures in about three to four months, depending on various conditions. Millennia of selective breeding have resulted in varieties that display a wide range of traits; e.g. suited for particular environments/latitudes, producing different ratios and compositions of terpenoids and cannabinoids (CBD, THC, CBG, CBC, CBN...etc.), fiber quality, oil/seed yield, etc. Hemp grown for fiber is planted closely, resulting in tall, slender plants with long fibers. The use of industrial hemp plant and its cultivation was commonplace until the 1900s when it was associated with its genetic sibling a.k.a. Drug-Type Cannabis species (which contain higher levels of psychoactive THC). Influential groups misconstrued hemp as a dangerous "drug", even though hemp is not a recreational drug and has the potential to be a sustainable and profitable crop for many farmers due to hemp's medical, structural and dietary uses. In the United States, the public's perception of hemp as marijuana has blocked hemp from becoming a useful crop and product," in spite of its vital importance prior to World War II. Ideally, according to Britain's Department for Environment, Food and Rural Affairs, the herb should be desiccated and harvested toward the end of flowering. This early cropping reduces the seed yield but improves the fiber yield and quality. The seeds are sown with grain drills or other conventional seeding equipment to a depth of . Greater seeding depths result in increased weed competition. Nitrogen should not be placed with the seed, but phosphate may be tolerated. The soil should have available 89 to 135 kg/ha of nitrogen, 46 kg/ha phosphorus, 67 kg/ha potassium, and 17 kg/ha sulfur. Organic fertilizers such as manure are one of the best methods of weed control. Cultivars In contrast to cannabis for medical use, varieties grown for fiber and seed have less than 0.3% THC and are unsuitable for producing hashish and marijuana. Present in industrial hemp, cannabidiol is a major constituent among some 560 compounds found in hemp. Cannabis sativa L. subsp. sativa var. sativa is the variety grown for industrial use, while C. sativa subsp. indica generally has poor fiber quality and female buds from this variety are primarily used for recreational and medicinal purposes. The major differences between the two types of plants are the appearance, and the amount of Δ9-tetrahydrocannabinol (THC) secreted in a resinous mixture by epidermal hairs called glandular trichomes, although they can also be distinguished genetically. Oilseed and fiber varieties of Cannabis approved for industrial hemp production produce only minute amounts of this psychoactive drug, not enough for any physical or psychological effects. Typically, hemp contains below 0.3% THC, while cultivars of Cannabis grown for medicinal or recreational use can contain anywhere from 2% to over 20%. Harvesting Smallholder plots are usually harvested by hand. The plants are cut at 2 to 3 cm above the soil and left on the ground to dry. Mechanical harvesting is now common, using specially adapted cutter-binders or simpler cutters. The cut hemp is laid in swathes to dry for up to four days. This was traditionally followed by retting, either water retting (the bundled hemp floats in water) or dew retting (the hemp remains on the ground and is affected by the moisture in dew and by molds and bacterial action). Pests Several arthropods can cause damage or injury to hemp plants, but the most serious species are associated with the Insecta class. The most problematic for outdoor crops are the voracious stem-boring caterpillars, which include the European corn borer, Ostrinia nubilalis, and the Eurasian hemp borer, Grapholita delineana. As the names imply, they target the stems reducing the structural integrity of the plant. Another lepidopteran, the corn earworm, Helicoverpa zea, is known to damage flowering parts and can be challenging to control. Other foliar pests, found in both indoor and outdoor crops, include the hemp russet mite, Aculops cannibicola, and cannabis aphid, Phorodon cannabis. They cause injury by reducing plant vigor because they feed on the phloem of the plant. Root feeders can be difficult to detect and control because of their below surface habitat. A number of beetle grubs and chafers are known to cause damage to hemp roots, including the flea beetle and Japanese beetle, Popillia Japonica. The rice root aphid, Rhopalosiphum rufiabdominale, has also been reported but primarily affects indoor growing facilities. Integrated pest management strategies should be employed to manage these pests with prevention and early detection being the foundation of a resilient program. Cultural and physical controls should be employed in conjunction with biological pest controls, chemical applications should only be used as a last resort. Diseases Hemp plants can be vulnerable to various pathogens, including bacteria, fungi, nematodes, viruses and other miscellaneous pathogens. Such diseases often lead to reduced fiber quality, stunted growth, and death of the plant. These diseases rarely affect the yield of a hemp field, so hemp production is not traditionally dependent on the use of pesticides. Environmental impact Hemp is considered by a 1998 study in Environmental Economics to be environmentally friendly due to a decrease of land use and other environmental impacts, indicating a possible decrease of ecological footprint in a US context compared to typical benchmarks. A 2010 study, however, that compared the production of paper specifically from hemp and eucalyptus concluded that "industrial hemp presents higher environmental impacts than eucalyptus paper"; however, the article also highlights that "there is scope for improving industrial hemp paper production". Hemp is also claimed to require few pesticides and no herbicides, and it has been called a carbon negative raw material. Results indicate that high yield of hemp may require high total nutrient levels (field plus fertilizer nutrients) similar to a high yielding wheat crop. A United Nations report endorses the versatility and sustainability of hemp and its productive potential in developing countries. Hemp uses a quarter of the water required by cotton, and absorbs more carbon dioxide than other crops and most trees. Producers The world-leading producer of hemp is China, which produces more than 70% of the world output. France ranks second with about a quarter of the world production. Smaller production occurs in the rest of Europe, Chile, and North Korea. Over 30 countries produce industrial hemp, including Australia, Austria, Canada, Chile, China, Denmark, Egypt, Finland, Germany, Greece, Hungary, India, Italy, Japan, Korea, Netherlands, New Zealand, Poland, Portugal, Romania, Russia, Slovenia, Spain, Sweden, Switzerland, Thailand, Turkey, the United Kingdom and Ukraine. The United Kingdom and Germany resumed commercial production in the 1990s. British production is mostly used as bedding for horses; other uses are under development. Companies in Canada, the UK, the United States, and Germany, among many others, process hemp seed into a growing range of food products and cosmetics; many traditional growing countries continue to produce textile-grade fiber. Air-dried stem yields in Ontario have from 1998 and onward ranged from 2.6 to 14.0 tons of dry, retted stalks per hectare (1–5.5 t/ac) at 12% moisture. Yields in Kent County, have averaged 8.75 t/ha (3.5 t/ac). Northern Ontario crops averaged 6.1 t/ha (2.5 t/ac) in 1998. Statistic for the European Union for 2008 to 2010 say that the average yield of hemp straw has varied between 6.3 and 7.3 ton per ha. Only a part of that is bast fiber. Around one ton of bast fiber and 2–3 tons of core material can be decorticated from 3–4 tons of good-quality, dry-retted straw. For an annual yield of this level is it in Ontario recommended to add nitrogen (N):70–110 kg/ha, phosphate (P2O5): up to 80 kg/ha and potash (K2O): 40–90 kg/ha. The average yield of dry hemp stalks in Europe was 6 ton/ha (2.4 ton/ac) in 2001 and 2002. FAO argue that an optimum yield of hemp fiber is more than 2 tons per ha, while average yields are around 650 kg/ha. Australia In the Australian states of Tasmania, Victoria, Queensland, Western Australia, New South Wales, and most recently, South Australia, the state governments have issued licenses to grow hemp for industrial use. The first to initiate modern research into the potential of cannabis was the state of Tasmania, which pioneered the licensing of hemp during the early 1990s. The state of Victoria was an early adopter in 1998, and has reissued the regulation in 2008. Queensland has allowed industrial production under license since 2002, where the issuance is controlled under the Drugs Misuse Act 1986. Western Australia enabled the cultivation, harvest and processing of hemp under its Industrial Hemp Act 2004, New South Wales now issues licenses under a law, the Hemp Industry Regulations Act 2008 (No 58), that came into effect as of 6 November 2008. Most recently, South Australia legalized industrial hemp under South Australia's Industrial Hemp Act 2017, which commenced on 12 November 2017. Canada Commercial production (including cultivation) of industrial hemp has been permitted in Canada since 1998 under licenses and authorization issued by Health Canada. In the early 1990s, industrial hemp agriculture in North America began with the Hemp Awareness Committee at the University of Manitoba. The Committee worked with the provincial government to get research and development assistance and was able to obtain test plot permits from the Canadian government. Their efforts led to the legalization of industrial hemp (hemp with only minute amounts of tetrahydrocannabinol) in Canada and the first harvest in 1998. In 2017, the cultivated area for hemp in the Prairie provinces include Saskatchewan with more than , Alberta with , and Manitoba with . Canadian hemp is cultivated mostly for its food value as hulled hemp seeds, hemp oils, and hemp protein powders, with only a small fraction devoted to production of hemp fiber used for construction and insulation. France France is Europe's biggest producer (and the world's second largest producer) with cultivated. 70–80% of the hemp fiber produced in 2003 was used for specialty pulp for cigarette papers and technical applications. About 15% was used in the automotive sector, and 5–6% was used for insulation mats. About 95% of hurds were used as animal bedding, while almost 5% was used in the building sector. In 2010–2011, a total of was cultivated with hemp in the EU, a decline compared with previous year. Russia and Ukraine From the 1950s to the 1980s, the Soviet Union was the world's largest producer of hemp ( in 1970). The main production areas were in Ukraine, the Kursk and Orel regions of Russia, and near the Polish border. Since its inception in 1931, the Hemp Breeding Department at the Institute of Bast Crops in Hlukhiv (Glukhov), Ukraine, has been one of the world's largest centers for developing new hemp varieties, focusing on improving fiber quality, per-hectare yields, and low THC content. After the collapse of the Soviet Union, the commercial cultivation of hemp declined sharply. However, at least an estimated 2.5 million acres of hemp grow wild in the Russian Far East and the Black Sea regions. United Kingdom In the United Kingdom, cultivation licenses are issued by the Home Office under the Misuse of Drugs Act 1971. When grown for nondrug purposes, hemp is referred to as industrial hemp, and a common product is fiber for use in a wide variety of products, as well as the seed for nutritional aspects and the oil. Feral hemp or ditch weed is usually a naturalized fiber or oilseed strain of Cannabis that has escaped from cultivation and is self-seeding. United States In October 2019, hemp became legal to grow in 46 U.S. states under federal law. As of 2019, 47 states have enacted legislation to make hemp legal to grow at the state level, with several states implementing medical provisions regarding the growing of plants specifically for non-psychoactive CBD. The 2018 Farm Bill, which incorporated the Hemp Farming Act of 2018, removed hemp as a Schedule I drug and instead made it an agricultural commodity. This legalized hemp at the federal level, which made it easier for hemp farmers to get production licenses, acquire loans, and receive federal crop insurance. NH 2014 N.H. Laws, Chap. 18, SD: HB 1008 (2020) S.D. Codified Laws Ann. §38-35-1 et seq. Authorizes the growth, production and transportation of hemp with a license, and directs the Department of Agriculture to submit a state plan to USDA. Requires a minimum of five contiguous outdoor acres for grower license applications, and requires any license applicants to submit to a state and federal criminal background investigation. Requires a transportation permit for any transporter traveling within or through the state and creates two types of industrial hemp transportation permits (grower licensee and general) provided by the Department of Public Safety. Creates the Hemp Regulatory Program Fund. The process to legalize hemp cultivation began in 2009, when Oregon began approving licenses for industrial hemp. Then, in 2013, after the legalization of marijuana, several farmers in Colorado planted and harvested several acres of hemp, bringing in the first hemp crop in the United States in over half a century. After that, the federal government created a Hemp Farming Pilot Program as a part of the Agricultural Act of 2014. This program allowed institutions of higher education and state agricultural departments to begin growing hemp without the consent of the Drug Enforcement Administration (DEA). Hemp production in Kentucky, formerly the United States' leading producer, resumed in 2014. Hemp production in North Carolina resumed in 2017, and in Washington State the same year. By the end of 2017, at least 34 U.S. states had industrial hemp programs. In 2018, New York began taking strides in industrial hemp production, along with hemp research pilot programs at Cornell University, Binghamton University and SUNY Morrisville. As of 2017, the hemp industry estimated that annual sales of hemp products were around $820 million annually; hemp-derived CBD have been the major force driving this growth. Despite this progress, hemp businesses in the US have had difficulties expanding as they have faced challenges in traditional marketing and sales approaches. According to a case study done by Forbes, hemp businesses and startups have had difficulty marketing and selling non-psychoactive hemp products, as majority of online advertising platforms and financial institutions do not distinguish between hemp and marijuana. History Gathered hemp fiber was used to make cloth long before agriculture, nine to fifty thousand years ago. It may also be one of the earliest plants to have been cultivated. An archeological site in the Oki Islands of Japan contained cannabis achenes from about 8000 BC, probably signifying use of the plant. Hemp use archaeologically dates back to the Neolithic Age in China, with hemp fiber imprints found on Yangshao culture pottery dating from the 5th millennium BC. The Chinese later used hemp to make clothes, shoes, ropes, and an early form of paper. The classical Greek historian Herodotus (ca. 480 BC) reported that the inhabitants of Scythia would often inhale the vapors of hemp-seed smoke, both as ritual and for their own pleasurable recreation. Textile expert Elizabeth Wayland Barber summarizes the historical evidence that Cannabis sativa, "grew and was known in the Neolithic period all across the northern latitudes, from Europe (Germany, Switzerland, Austria, Romania, Ukraine) to East Asia (Tibet and China)," but, "textile use of Cannabis sativa does not surface for certain in the West until relatively late, namely the Iron Age." "I strongly suspect, however, that what catapulted hemp to sudden fame and fortune as a cultigen and caused it to spread rapidly westwards in the first millennium B.C. was the spread of the habit of pot-smoking from somewhere in south-central Asia, where the drug-bearing variety of the plant originally occurred. The linguistic evidence strongly supports this theory, both as to time and direction of spread and as to cause." Jews living in Palestine in the 2nd century were familiar with the cultivation of hemp, as witnessed by a reference to it in the Mishna (Kil'ayim 2:5) as a variety of plant, along with arum, that sometimes takes as many as three years to grow from a seedling. In late medieval Holy Roman Empire (Germany) and Italy, hemp was employed in cooked dishes, as filling in pies and tortes, or boiled in a soup. Hemp in later Europe was mainly cultivated for its fibers and was used for ropes on many ships, including those of Christopher Columbus. The use of hemp as a cloth was centered largely in the countryside, with higher quality textiles being available in the towns. The Spaniards brought hemp to the Americas and cultivated it in Chile starting about 1545. Similar attempts were made in Peru, Colombia, and Mexico, but only in Chile did the crop find success. In July 1605, Samuel Champlain reported the use of grass and hemp clothing by the (Wampanoag) people of Cape Cod and the (Nauset) people of Plymouth Bay told him they harvested hemp in their region where it grew wild to a height of 4 to 5 ft. In May 1607, "hempe" was among the crops Gabriel Archer observed being cultivated by the natives at the main Powhatan village, where Richmond, Virginia, is now situated; and in 1613, Samuell Argall reported wild hemp "better than that in England" growing along the shores of the upper Potomac. As early as 1619, the first Virginia House of Burgesses passed an Act requiring all planters in Virginia to sow "both English and Indian" hemp on their plantations. The Puritans are first known to have cultivated hemp in New England in 1645. United States George Washington pushed for the growth of hemp as it was a cash crop commonly used to make rope and fabric. In May 1765 he noted in his diary about the sowing of seeds each day until mid-April. Then he recounts the harvest in October which he grew 27 bushels that year. It is sometimes supposed that an excerpt from Washington's diary, which reads "Began to the Male from the Female hemp at Do.&—rather too late" is evidence that he was trying to grow female plants for the THC found in the flowers. However, the editorial remark accompanying the diary states that "This may arise from their [the male] being coarser, and the stalks larger" In subsequent days, he describes soaking the hemp (to make the fibers usable) and harvesting the seeds, suggesting that he was growing hemp for industrial purposes, not recreational. George Washington also imported the Indian hemp plant from Asia, which was used for fiber and, by some growers, for intoxicating resin production. In a 1796 letter to William Pearce who managed the plants for him, Washington says, "What was done with the Indian Hemp plant from last summer? It ought, all of it, to be sown again; that not only a stock of seed sufficient for my own purposes might have been raised, but to have disseminated seed to others; as it is more valuable than common hemp." Other presidents known to have farmed hemp for alternative purposes include Thomas Jefferson, James Madison, James Monroe, Andrew Jackson, Zachary Taylor, and Franklin Pierce. Historically, hemp production had made up a significant portion of antebellum Kentucky's economy. Before the American Civil War, many slaves worked on plantations producing hemp. In 1937, the Marihuana Tax Act of 1937 was passed in the United States, levying a tax on anyone who dealt commercially in cannabis, hemp, or marijuana. The passing of the Act to destroy the U.S. hemp industry has been reputed to involve businessmen Andrew Mellon, Randolph Hearst and the Du Pont family. One claim is that Hearst believed that his extensive timber holdings were threatened by the invention of the decorticator that he feared would allow hemp to become a cheap substitute for the paper pulp used for newspaper. Historical research indicates this fear was unfounded because improvements of the decorticators in the 1930s – machines that separated the fibers from the hemp stem – could not make hemp fiber a cheaper substitute for fibers from other sources. Further, decorticators did not perform satisfactorily in commercial production. Another claim is that Mellon, Secretary of the Treasury and the wealthiest man in America at that time, had invested heavily in DuPont's new synthetic fiber, nylon, and believed that the replacement of the traditional resource, hemp, was integral to the new product's success. DuPont and many industrial historians dispute a link between nylon and hemp, nylon became immediately a scarce commodity. Nylon had characteristics that could be used for toothbrushes (sold from 1938) and very thin nylon fiber could compete with silk and rayon in various textiles normally not produced from hemp fiber, such as very thin stockings for women. While the Marijuana Tax Act of 1937 had just been signed into law, the United States Department of Agriculture lifted the tax on hemp cultivation during WWII. Before WWII, the U.S. Navy used Jute and Manila Hemp from the Philippines and Indonesia for the cordage on their ships. During the war, Japan cut off those supply lines. America was forced to turn inward and revitalize the cultivation of Hemp on U.S. soils. Hemp was used extensively by the United States during World War II to make uniforms, canvas, and rope. Much of the hemp used was cultivated in Kentucky and the Midwest. During World War II, the U.S. produced a short 1942 film, Hemp for Victory, promoting hemp as a necessary crop to win the war. By the 1980s the film was largely forgotten, and the U.S. government even denied its existence. The film, and the important historical role of hemp in U.S. agriculture and commerce was brought to light by hemp activist Jack Herer in the book The Emperor Wears No Clothes. U.S. farmers participated in the campaign to increase U.S. hemp production to 36,000 acres in 1942. This increase amounted to more than 20 times the production in 1941 before the war effort. In the United States, Executive Order 12919 (1994) identified hemp as a strategic national product that should be stockpiled. Historical cultivation Hemp has been grown for millennia in Asia and the Middle East for its fiber. Commercial production of hemp in the West took off in the eighteenth century, but was grown in the sixteenth century in eastern England. Because of colonial and naval expansion of the era, economies needed large quantities of hemp for rope and oakum. In the early 1940s, world production of hemp fiber ranged from 250,000 to 350,000 metric tons, Russia was the biggest producer. In Western Europe, the cultivation of hemp was not legally banned by the 1930s, but the commercial cultivation stopped by then, due to decreased demand compared to increasingly popular artificial fibers. Speculation about the potential for commercial cultivation of hemp in large quantities has been criticized due to successful competition from other fibers for many products. The world production of hemp fiber fell from over 300,000 metric tons 1961 to about 75,000 metric tons in the early 1990s and has after that been stable at that level. Japan In Japan, hemp was historically used as paper and a fiber crop. There is archaeological evidence cannabis was used for clothing and the seeds were eaten in Japan back to the Jōmon period (10,000 to 300 BC). Many Kimono designs portray hemp, or asa (), as a beautiful plant. In 1948, marijuana was restricted as a narcotic drug. The ban on marijuana imposed by the United States authorities was alien to Japanese culture, as the drug had never been widely used in Japan before. Though these laws against marijuana are some of the world's strictest, allowing five years imprisonment for possession of the drug, they exempt hemp growers, whose crop is used to make robes for Buddhist monks and loincloths for Sumo wrestlers. Because marijuana use in Japan has doubled in the past decade, these exemptions have recently been called into question. Portugal The cultivation of hemp in Portuguese lands began around the fourteenth century. The raw material was used for the preparation of rope and plugs for the Portuguese ships. Portugal also utilized its colonies to support its hemp supply, including in certain parts of Brazil. In order to recover the ailing Portuguese naval fleet after the Restoration of Independence in 1640, King John IV put a renewed emphasis on the growing of hemp. He ordered the creation of the Royal Linen and Hemp Factory in the town of Torre de Moncorvo to increase production and support the effort. In 1971, the cultivation of hemp became illegal, and the production was substantially reduced. Because of EU regulations 1308–70, 619/71 and 1164–89, this law was revoked (for some certified seed varieties).	Biology and health sciences	Rosales	Plants
963403	https://en.wikipedia.org/wiki/Messier%2049	Messier 49	Messier 49 (also known as M49 or NGC 4472) is a giant elliptical galaxy about away in the equatorial constellation of Virgo. This galaxy was discovered by astronomer Charles Messier in 1777. As an elliptical galaxy, Messier 49 has the physical form of a radio galaxy, but it only has the radio emission of a normal galaxy. From the detected radio emission, the core region has roughly 1053 erg (1046 J or 1022 YJ) of synchrotron energy. The nucleus of this galaxy is emitting X-rays, suggesting the likely presence of a supermassive black hole with an estimated mass of , or 565 million times the mass of the Sun (). X-ray emissions shows a structure to the north of Messier 49 that resembles a bow shock. To the southwest of the core, the luminous outline of the galaxy can be traced out to a distance of 260 kpc. This galaxy has many globular clusters: estimated to be about 5,900. This is far more than the roughly 200 orbiting the Milky Way, but dwarfed by the 13,450 orbiting the supergiant elliptical galaxy Messier 87. On average, the globular clusters of M49 are about 10 billion years old. Between 2000 and 2009, strong evidence for a stellar mass black hole was discovered in one. A second candidate was announced in 2011. Messier 49 was the first member of the Virgo Cluster of galaxies to be discovered. It is the most luminous member of that cluster and more luminous than any galaxy closer to the Earth. This galaxy forms part of the smaller Virgo B subcluster 4.5° away from the dynamic center of the Virgo Cluster, centered on Messier 87. Messier 49 is gravitationally interacting with the dwarf irregular galaxy UGC 7636. The dwarf shows a trail of debris spanning roughly 1 × 5 arcminutes, which corresponds to a physical dimension of . One supernova has been observed in M49: SN 1969Q (type unknown, mag. 13) was discovered by Evans on 12 June 1969. [Note: some sources incorrectly report the discovery date as 1 June 1969.]	Physical sciences	Notable galaxies	Astronomy
963967	https://en.wikipedia.org/wiki/Fireplace	Fireplace	A fireplace or hearth is a structure made of brick, stone or metal designed to contain a fire. Fireplaces are used for the relaxing ambiance they create and for heating a room. Modern fireplaces vary in heat efficiency, depending on the design. Historically, they were used for heating a dwelling, cooking, and heating water for laundry and domestic uses. A fire is contained in a firebox or fire pit; a chimney or other flue allows exhaust gas to escape. A fireplace may have the following: a foundation, a hearth, a firebox, a mantel, a chimney crane (used in kitchen and laundry fireplaces), a grate, a lintel, a lintel bar, an overmantel, a damper, a smoke chamber, a throat, a flue, and a chimney filter or afterburner. On the exterior, there is often a corbelled brick crown, in which the projecting courses of brick act as a drip course to keep rainwater from running down the exterior walls. A cap, hood, or shroud serves to keep rainwater out of the exterior of the chimney; rain in the chimney is a much greater problem in chimneys lined with impervious flue tiles or metal liners than with the traditional masonry chimney, which soaks up all but the most violent rain. Some chimneys have a spark arrestor incorporated into the crown or cap. Organizations like the United States Environmental Protection Agency (EPA) and the Washington State Department of Ecology warn that, according to various studies, fireplaces can pose health risks. The EPA writes "Smoke may smell good, but it's not good for you." Types of fireplaces Manufactured fireplaces are made with sheet metal or glass fire boxes. Electric fireplaces can be built-in replacements for wood or gas or retrofit with log inserts or electric fireboxes. A few types are wall mounted electric fireplaces, electric fireplace stoves, electric mantel fireplaces, and fixed or free standing electric fireplaces. Masonry and prefabricated fireplaces can be fueled by: Wood fuel or firewood and other biomass Charcoal (carbonized biomass) Coal of various grades Coke (carbonized coal) Smokeless fuel of several types Flammable gases: propane, butane, and methane (natural gas is mostly methane, liquefied petroleum gas mostly propane) Ethanol (a liquid alcohol, also sold in gels) Ventless fireplaces (duct free/room-venting fireplaces) are fueled by either gel, liquid propane, bottled gas or natural gas. In the United States, some states and local counties have laws restricting these types of fireplaces. They must be properly sized to the area to be heated. There are also air quality control issues due to the amount of moisture they release into the room air, and an oxygen sensor and a carbon monoxide detector are safety essentials. Direct vent fireplaces are fueled by either liquid propane or natural gas. They are completely sealed from the area that is heated, and vent all exhaust gasses to the exterior of the structure. Chimney and flue types: Masonry (brick or stone fireplaces and chimneys) with or without tile-lined flue. Reinforced concrete chimneys. Fundamental design flaws bankrupted the US manufacturers and made the design obsolete. These chimneys often show vertical cracks on the exterior. Metal-lined flue: Double- or triple-walled metal pipe running up inside a new or existing wood-framed or masonry chase. Newly constructed flues may feature a chase cover, a cap, and a spark arrestor at the top to keep small animals out and to prevent sparks from being broadcast into the atmosphere. All gas fireplaces require trained gas service members to carry out installations. Accessories A wide range of accessories are used with fireplaces, which range between countries, regions, and historical periods. For the interior, common in recent Western cultures include grates, fireguards, log boxes, andirons and pellet baskets, all of which cradle fuel and accelerate combustion. A grate (or fire grate) is a frame, usually of iron bars, to retain fuel for a fire. Heavy metal firebacks are sometimes used to capture and re-radiate heat, to protect the back of the fireplace, and as decoration. Fenders are low metal frames set in front of the fireplace to contain embers, soot and ash. For fireplace tending, tools include pokers, bellows, tongs, shovels, brushes and tool stands. Other wider accessories can include log baskets, companion sets, coal buckets, cabinet accessories and more. History Ancient fire pits were sometimes built in the ground, within caves, or in the center of a hut or dwelling. Evidence of prehistoric, man-made fires exists on all six inhabited continents. The disadvantage of early indoor fire pits was that they produced toxic and/or irritating smoke inside the dwelling. Fire pits developed into raised hearths in buildings, but venting smoke depended on open windows or holes in roofs. The medieval great hall typically had a centrally located hearth, where an open fire burned with the smoke rising to the vent in the roof. Louvers were developed during the Middle Ages to allow the roof vents to be covered so rain and snow would not enter. Also during the Middle Ages, smoke canopies were invented to prevent smoke from spreading through a room and vent it out through a wall or roof. These could be placed against stone walls, instead of taking up the middle of the room, and this allowed smaller rooms to be heated. Chimneys were invented in northern Europe in the 11th or 12th century and largely fixed the problem of smoke, more reliably venting it outside. They made it possible to give the fireplace a draft, and also made it possible to put fireplaces in multiple rooms in buildings conveniently. They did not come into general use immediately, however, as they were expensive to build and maintain. In 1678, Prince Rupert, nephew of Charles I, raised the grate of the fireplace, improving the airflow and venting system. The 18th century saw two important developments in the history of fireplaces. Benjamin Franklin developed a convection chamber for the fireplace that greatly improved the efficiency of fireplaces and wood-burning stoves. He also improved the airflow by pulling air from a basement and venting out a longer area at the top. In the later 18th century, Count Rumford designed a fireplace with a tall, shallow firebox that was better at drawing the smoke up and out of the building. The shallow design also improved greatly the amount of heat transfer projected into the room. The Aesthetic movement of the 1870s and 1880s favoured a more traditional look based on stone, with simple designs and limited ornamentation. In the 1890s, the Aesthetic movement gave way to the Arts and Crafts movement, which still emphasized quality stone and practical features. Stone fireplaces at this time were a symbol of prosperity, as to some degree they remain today. Evolution of fireplace design Over time, the purpose of fireplaces has changed from one of necessity to one of visual interest. Early ones were more fire pits than modern fireplaces. They were used for warmth on cold days and nights, as well as for cooking. They also served as a gathering place within the home. These fire pits were usually centered within a room, allowing more people to gather around it. Many flaws were found in early fireplace designs. Along with the Industrial Revolution, came large-scale housing developments, necessitating a standardization of fireplaces. The most renowned fireplace designers of this time were the Adam Brothers: John Adam, Robert Adam, and James Adam. They perfected a style of fireplace design that was used for generations. It was smaller, more brightly lit, with an emphasis on the quality of the materials used in their construction, instead of their size. By the 1800s, most new fireplaces were made up of two parts, the surround and the insert. The surround consisted of the mantelpiece and side supports, usually in wood, marble or granite. The insert was where the fire burned, and was constructed of cast iron often backed with decorative tiles. As well as providing heat, the fireplaces of the Victorian era were thought to add a cosy ambiance to homes. In the US state of Wisconsin, some elementary classrooms would contain decorated fireplaces to ease children's transition from home to school. Heating efficiency Some fireplace units incorporate a blower, which transfers more of the fireplace's heat to the air via convection, resulting in a more evenly heated space and a lower heating load. Fireplace efficiency can also be increased with the use of a fireback, a piece of metal that sits behind the fire and reflects heat back into the room. Firebacks are traditionally made from cast iron, but are also made from stainless steel. Most older fireplaces have a relatively low efficiency rating. Standard, modern, wood-burning masonry fireplaces though have an efficiency rating of at least 80% (legal minimum requirement, for example, in Salzburg, Austria). To improve efficiency, fireplaces can also be modified by inserting special heavy fireboxes designed to burn much cleaner and can reach efficiencies as high as 80% in heating the air. These modified fireplaces are often equipped with a large fire window, enabling an efficient heating process in two phases. During the first phase the initial heat is provided through a large glass window while the fire is burning. During this time the structure, built of refractory bricks, absorbs the heat. This heat is then evenly radiated for many hours during the second phase. Masonry fireplaces without a glass fire window only provide heat radiated from its surface. Depending on the outside temperature, 1 to 2 daily firings are sufficient to ensure a constant room temperature. Health effects Wood A literature review published in the Journal of Toxicology and Environmental Health concludes that there are a wide variety of health risks posed by residential wood combustion. It states: The Washington State Department of Ecology also published a booklet explaining why wood smoke can be dangerous. It explains that human lung and respiratory systems are unable to filter particulates emitted by wood combustion, which penetrate deeply into the lungs. For months, carcinogens can continue to cause changes and structural damage within the respiratory system. Young children, seniors, pregnant women, smokers and individuals with respiratory diseases are most vulnerable. Wood smoke can cause disease and even death in children, because it is associated with lower respiratory tract infections. Home fireplaces have caused fatal carbon monoxide poisoning. Gases and ethanol Propane, butane, and methane are all flammable gases used in fireplaces (natural gas is mostly methane, liquefied petroleum gas mostly propane). Gases can act as asphyxiant gases or cause gas explosions if they are allowed to accumulate unburned. Ethanol (a liquid, also sold in gels) fires can also cause severe burns. Burning hydrocarbons can decrease indoor air quality. Emissions include airborne particulate matter (such as black carbon) and gases like nitrogen oxide. These harm health: they weaken the immune system, and increase infections, blood pressure, cardiovascular diseases, and insulin resistance. Some forms of fuel are more harmful than others. Burning hydrocarbon fuels incompletely can produce carbon monoxide, which is highly poisonous and can cause death and long-term neurological disorders. Environmental effects Burning any hydrocarbon fuel releases carbon dioxide and water vapor. Other emissions, such as nitrogen oxides and sulfur oxides, can be harmful to the environment. Glossary Several of these terms may be compounded with chimney or fireplace such as chimney-back. Andiron—Either one of two horizontal metal bars resting on short legs intended to support firewood in a hearth. Arch—An arched top of the fireplace opening. Ash dump—An opening in a hearth to sweep ashes for later removal from the ash pit. Back (fireback)—The inside, rear wall of the fireplace of masonry or metal that reflects heat into the room. Brick trimmer—A brick arch supporting a hearth or shielding a joist in front of a fireplace. Chimney breast—The part of the chimney which projects into a room to accommodate a fireplace. Crane—Metal arms mounted on pintles, which swing and hold pots above a fire. Damper—A metal door to close a flue when a fireplace is not in use. Flue—The passageway in the chimney. Hearth—The floor of a fireplace. The part of a hearth which projects into a room may be called the front or outer hearth. Hearthstone—A large stone or other materials used as the hearth material. Insert—The fireplace insert is a device inserted into an existing masonry or prefabricated wood fireplace. Jamb—The side of a fireplace opening. Mantel—Either the shelf above a fireplace or the structure to support masonry above a fireplace Smoke shelf—A shelf below the smoke chamber and behind the damper. It collects debris and water falling down the flue. Throat (waist)—The narrow area above a fireplace usually where the damper is located. Wing—The sides of a fireplace above the opening near the throat.	Technology	Other components	null
964428	https://en.wikipedia.org/wiki/Weighing%20scale	Weighing scale	A scale or balance is a device used to measure weight or mass. These are also known as mass scales, weight scales, mass balances, massometers, and weight balances. The traditional scale consists of two plates or bowls suspended at equal distances from a fulcrum. One plate holds an object of unknown mass (or weight), while objects of known mass or weight, called weights, are added to the other plate until mechanical equilibrium is achieved and the plates level off, which happens when the masses on the two plates are equal. The perfect scale rests at neutral. A spring scale will make use of a spring of known stiffness to determine mass (or weight). Suspending a certain mass will extend the spring by a certain amount depending on the spring's stiffness (or spring constant). The heavier the object, the more the spring stretches, as described in Hooke's law. Other types of scales making use of different physical principles also exist. Some scales can be calibrated to read in units of force (weight) such as newtons instead of units of mass such as kilograms. Scales and balances are widely used in commerce, as many products are sold and packaged by mass. Pan balance History The balance scale is such a simple device that its usage likely far predates the evidence. What has allowed archaeologists to link artifacts to weighing scales are the stones for determining absolute mass. The balance scale itself was probably used to determine relative mass long before absolute mass. The oldest attested evidence for the existence of weighing scales dates to the Fourth Dynasty of Egypt, with Deben (unit) balance weights, from the reign of Sneferu (c. 2600 BC) excavated, though earlier usage has been proposed. Carved stones bearing marks denoting mass and the Egyptian hieroglyphic symbol for gold have been discovered, which suggests that Egyptian merchants had been using an established system of mass measurement to catalog gold shipments or gold mine yields. Although no actual scales from this era have survived, many sets of weighing stones as well as murals depicting the use of balance scales suggest widespread usage. Examples, dating , have also been found in the Indus River valley. Uniform, polished stone cubes discovered in early settlements were probably used as mass-setting stones in balance scales. Although the cubes bear no markings, their masses are multiples of a common denominator. The cubes are made of many different kinds of stones with varying densities. Clearly their mass, not their size or other characteristics, was a factor in sculpting these cubes. In China, the earliest weighing balance excavated was from a tomb of the State of Chu of the Chinese Warring States Period dating back to the 3rd to 4th century BC in Mount Zuojiagong near Changsha, Hunan. The balance was made of wood and used bronze masses. Variations on the balance scale, including devices like the cheap and inaccurate bismar (unequal-armed scales), began to see common usage by c. 400 BC by many small merchants and their customers. A plethora of scale varieties each boasting advantages and improvements over one another appear throughout recorded history, with such great inventors as Leonardo da Vinci lending a personal hand in their development. Even with all the advances in weighing scale design and development, all scales until the seventeenth century AD were variations on the balance scale. The standardization of the weights used – and ensuring traders used the correct weights – was a considerable preoccupation of governments throughout this time. The original form of a balance consisted of a beam with a fulcrum at its center. For highest accuracy, the fulcrum would consist of a sharp V-shaped pivot seated in a shallower V-shaped bearing. To determine the mass of the object, a combination of reference masses was hung on one end of the beam while the object of unknown mass was hung on the other end (see balance and steelyard balance). For high precision work, such as empirical chemistry, the center beam balance is still one of the most accurate technologies available, and is commonly used for calibrating test masses. However, bronze fragments discovered in central Germany and Italy had been used during the Bronze Age as an early form of currency. In the same time period, merchants had used standard weights of equivalent value between 8 and 10.5 grams from Great Britain to Mesopotamia. Mechanical balances The balance (also balance scale, beam balance and laboratory balance) was the first mass measuring instrument invented. In its traditional form, it consists of a pivoted horizontal lever with arms of equal lengththe beam or tron and a weighing pan suspended from each arm (hence the plural name "scales for a weighing instrument). The unknown mass is placed in one pan and standard masses are added to the other pan until the beam is as close to equilibrium as possible. In precision balances, a more accurate determination of the mass is given by the position of a sliding mass moved along a graduated scale. A decimal balance uses the lever in which the arm for weights is 10 times longer than the arm for weighted objects, so that much lighter weights may be used to weigh heavy object. Similarly a centesimal balance uses arms in ratio 1:100. Unlike spring-based scales, balances are used for the precision measurement of mass as their accuracy is not affected by variations in the local gravitational field. (On Earth, for example, these can amount to ±0.5% between locations.) A change in the strength of the gravitational field caused by moving the balance does not change the measured mass, because the moments of force on either side of the center balanced beam are affected equally. A center beam balance will render an accurate measurement of mass at any location experiencing a constant gravity or acceleration. Very precise measurements are achieved by ensuring that the balance's fulcrum is essentially friction-free (a knife edge is the traditional solution), by attaching a pointer to the beam which amplifies any deviation from a balance position; and finally by using the lever principle, which allows fractional masses to be applied by movement of a small mass along the measuring arm of the beam, as described above. For greatest accuracy, there needs to be an allowance for the buoyancy in air, whose effect depends on the densities of the masses involved. To reduce the need for large reference masses, an off-center beam can be used. A balance with an off-center beam can be almost as accurate as a scale with a center beam, but the off-center beam requires special reference masses and cannot be intrinsically checked for accuracy by simply swapping the contents of the pans as a center-beam balance can. To reduce the need for small graduated reference masses, a sliding weight called a poise can be installed so that it can be positioned along a calibrated scale. A poise adds further intricacies to the calibration procedure, since the exact mass of the poise must be adjusted to the exact lever ratio of the beam. For greater convenience in placing large and awkward loads, a platform can be floated on a cantilever beam system which brings the proportional force to a noseiron bearing; this pulls on a stilyard rod to transmit the reduced force to a conveniently sized beam. One still sees this design in portable beam balances of 500 kg capacity which are commonly used in harsh environments without electricity, as well as in the lighter duty mechanical bathroom scale (which actually uses a spring scale, internally). The additional pivots and bearings all reduce the accuracy and complicate calibration; the float system must be corrected for corner errors before the span is corrected by adjusting the balance beam and poise. Roberval balance In 1669 the Frenchman Gilles Personne de Roberval presented a new kind of balance scale to the French Academy of Sciences. This scale consisted of a pair of vertical columns separated by a pair of equal-length arms and pivoting in the center of each arm from a central vertical column, creating a parallelogram. From the side of each vertical column a peg extended. To the amazement of observers, no matter where Roberval hung two equal weight along the peg, the scale still balanced. In this sense, the scale was revolutionary: it evolved into the more-commonly encountered form consisting of two pans placed on vertical column located above the fulcrum and the parallelogram below them. The advantage of the Roberval design is that no matter where equal weights are placed in the pans, the scale will still balance. Further developments have included a "gear balance" in which the parallelogram is replaced by any odd number of interlocking gears greater than one, with alternating gears of the same size and with the central gear fixed to a stand and the outside gears fixed to pans, as well as the "sprocket gear balance" consisting of a bicycle-type chain looped around an odd number of sprockets with the central one fixed and the outermost two free to pivot and attached to a pan. Because it has more moving joints which add friction, the Roberval balance is consistently less accurate than the traditional beam balance, but for many purposes this is compensated for by its usability. Torsion balance The torsion balance is one of the most mechanically accurate of analog balances. Pharmacy schools still teach how to use torsion balances in the U.S. It utilizes pans like a traditional balance that lie on top of a mechanical chamber which bases measurements on the amount of twisting of a wire or fiber inside the chamber. The scale must still use a calibration weight to compare against, and can weigh objects greater than 120 mg and come within a margin of error +/- 7 mg. Many microbalances and ultra-microbalances that weigh fractional gram values are torsion balances. A common fiber type is quartz crystal. Electronic devices Microbalance A microbalance (also called an ultramicrobalance, or nanobalance) is an instrument capable of making precise measurements of the mass of objects of relatively small mass: on the order of a million parts of a gram and below. Analytical balance An analytical balance is a class of balance designed to measure small mass in the sub-milligram range. The measuring pan of an analytical balance (0.1 mg or better) is inside a transparent enclosure with doors so that dust does not collect and so any air currents in the room do not affect the balance's operation. This enclosure is often called a draft shield. The use of a mechanically vented balance safety enclosure, which has uniquely designed acrylic airfoils, allows a smooth turbulence-free airflow that prevents balance fluctuation and the measure of mass down to 1 μg without fluctuations or loss of product. Also, the sample must be at room temperature to prevent natural convection from forming air currents inside the enclosure from causing an error in reading. Single-pan mechanical substitution balances maintain consistent response throughout the useful capacity, which is achieved by maintaining a constant load on the balance beam and thus the fulcrum by subtracting mass on the same side of the beam to which the sample is added. Electronic analytical scales measure the force needed to counter the mass being measured rather than using actual masses. As such they must have calibration adjustments made to compensate for gravitational differences. They use an electromagnet to generate a force to counter the sample being measured and output the result by measuring the force needed to achieve balance. Such a measurement device is called an electromagnetic force restoration sensor. Pendulum balance scales Pendulum type scales do not use springs. These designs use pendulums and operate as a balance that is unaffected by differences in gravity. An example of application of this design are scales made by the Toledo Scale Company. Programmable scales A programmable scale has a programmable logic controller in it, allowing it to be programmed for various applications such as batching, labeling, filling (with check weight function), truck scales, and more. Another important function is counting, e. g. used to count small parts in larger quantities during the annual stock taking. Counting scales (which can also do just weighing) can range from mg to tonnes. Symbolism The scales (specifically, a two-pan, beam balance) are one of the traditional symbols of justice, as wielded by statues of Lady Justice. This corresponds to the use in a metaphor of matters being "held in the balance". It has its origins in ancient Egypt. Scales also are widely used as a symbol of finance, commerce, or trade, in which they have played a traditional, vital role since ancient times. For instance, balance scales are depicted in the seal of the U.S. Department of the Treasury and the Federal Trade Commission. Scales are also the symbol for the astrological sign Libra. Scales (specifically, a two-pan, beam balance in a state of equal balance) are the traditional symbol of Pyrrhonism indicating the equal balance of arguments used in inducing epoche. Force-measuring (weight) scales History Although records dating to the 1700s refer to spring scales for measuring mass, the earliest design for such a device dates to 1770 and credits Richard Salter, an early scale-maker. Spring scales came into wide usage in the United Kingdom after 1840 when R. W. Winfield developed the candlestick scale for weighing letters and packages, required after the introduction of the Uniform Penny Post. Postal workers could work more quickly with spring scales than balance scales because they could be read instantaneously and did not have to be carefully balanced with each measurement. By the 1940s, various electronic devices were being attached to these designs to make readings more accurate. Load cells – transducers that convert force to an electrical signal – have their beginnings as early as the late nineteenth century, but it was not until the late twentieth century that their widespread usage became economically and technologically viable. Mechanical scales A mechanical scale or balance is used to describe a weighing device that is used to measure the mass, force exertion, tension, and resistance of an object without the need of a power supply. Types of mechanical scales include decimal balances, spring scales, hanging scales, triple beam balances, and force gauges. Spring scales A spring scale measures mass by reporting the distance that a spring deflects under a load. This contrasts to a balance, which compares the torque on the arm due to a sample weight to the torque on the arm due to a standard reference mass using a horizontal lever. Spring scales measure force, which is the tension force of constraint acting on an object, opposing the local force of gravity. They are usually calibrated so that measured force translates to mass at earth's gravity. The object to be weighed can be simply hung from the spring or set on a pivot and bearing platform. In a spring scale, the spring either stretches (as in a hanging scale in the produce department of a grocery store) or compresses (as in a simple bathroom scale). By Hooke's law, every spring has a proportionality constant that relates how hard it is pulled to how far it stretches. Weighing scales use a spring with a known spring constant (see Hooke's law) and measure the displacement of the spring by any variety of mechanisms to produce an estimate of the gravitational force applied by the object. Rack and pinion mechanisms are often used to convert the linear spring motion to a dial reading. Spring scales have two sources of error that balances do not: the measured mass varies with the strength of the local gravitational force (by as much as 0.5% at different locations on Earth), and the elasticity of the measurement spring can vary slightly with temperature. With proper manufacturing and setup, however, spring scales can be rated as legal for commerce. To remove the temperature error, a commerce-legal spring scale must either have temperature-compensated springs or be used at a fairly constant temperature. To eliminate the effect of gravity variations, a commerce-legal spring scale must be calibrated where it is used. Hydraulic or pneumatic scale It is also common in high-capacity applications such as crane scales to use hydraulic force to sense mass. The test force is applied to a piston or diaphragm and transmitted through hydraulic lines to a dial indicator based on a Bourdon tube or electronic sensor. Domestic Weighing Scale Electronic digital scales display weight as a number, usually on a liquid crystal display (LCD). They are versatile because they may perform calculations on the measurement and transmit it to other digital devices. On a digital scale, the force of the weight causes a spring to deform, and the amount of deformation is measured by one or more transducers called strain gauges. A strain gauge is a conductor whose electrical resistance changes when its length changes. Strain gauges have limited capacity and larger digital scales may use a hydraulic transducer called a load cell instead. A voltage is applied to the device, and the weight causes the current through it to change. The current is converted to a digital number by an analog-to-digital converter, translated by digital logic to the correct units, and displayed on the display. Usually, the device is run by a microprocessor chip. Digital bathroom scale A digital bathroom scale is a scale on the floor which a person stands on. The weight is shown on an LED or LCD display. The digital electronics may do more than just display weight, it may calculate body fat, BMI, lean mass, muscle mass, and water ratio. Some modern bathroom scales are wirelessly or cellularly connected and have features like smartphone integration, cloud storage, and fitness tracking. They are usually powered by a button cell, or battery of AA or AAA size. Digital kitchen scale Digital kitchen scales are used for weighing food in a kitchen during cooking. These are usually lightweight and compact. Strain gauge scale In electronic versions of spring scales, the deflection of a beam supporting the unknown mass is measured using a strain gauge, which is a length-sensitive electrical resistance. The capacity of such devices is only limited by the resistance of the beam to deflection. The results from several supporting locations may be added electronically, so this technique is suitable for determining the mass of very heavy objects, such as trucks and rail cars, and is used in a modern weighbridge. Supermarket and other retail scale These scales are used in the modern bakery, grocery, delicatessen, seafood, meat, produce and other perishable goods departments. Supermarket scales can print labels and receipts, mark mass and count, unit price, total price and in some cases tare. Some modern supermarket scales print an RFID tag that can be used to track the item for tampering or returns. In most cases, these types of scales have a sealed calibration so that the reading on the display is correct and cannot be tampered with. In the US, the scales are certified by the National Type Evaluation Program (NTEP), in South Africa by the South African Bureau of Standards, in Australia, they are certified by the National Measurement Institute (NMI) and in the UK by the International Organization of Legal Metrology. Industrial weighing scale An industrial weighing scale is a device that measures the weight or mass of objects in various industries. It can range from small bench scales to large weighbridges, and it can have different features and capacities. Industrial weighing scales are used for quality control, inventory management, and trade purposes. There are many kinds of industrial weighing scales that are used for different purposes and applications. Some of the common types are: Weighbridges : A large scale that can weigh trucks, lorries, containers, and other heavy-duty vehicles. They are used in industries like manufacturing, shipping, mining, agriculture, etc Container Stacker Scale : A container stacker scale is a specialized weighing system designed for accurately measuring the weight of shipping containers. It is typically integrated into the equipment used for loading and unloading containers, such as container handlers or stacker cranes. Container stacker scales provide real-time weight measurements, allowing logistics professionals to ensure that each container is loaded within the specified weight limits. Container stacker scales are used in industries like ports, shipping, and logistics Forklift scale : A forklift scale is a weighing system that is built into a forklift truck. It allows for the weighing of loads while they are being lifted and transported by the forklift. This eliminates the need for separate weighing operations and reduces the time and labor required for material handling operations. Forklift scales are used in various industries, such as manufacturing, logistics, and shipping. Material Handler Scale : A Material Handler Scale is a weighing system that is integrated into a material handler machine, such as a grapple or a magnet. It allows for the accurate and efficient weighing of materials while they are being moved, unloaded, or loaded. A Material Handler Scale can be used in various industries, such as scrap, recycling, waste, and port and harbor. A Material Handler Scale can also transfer the weighing information to a cloud service or an ERP system for real-time monitoring and management of material flow. A pallet jack scale is a device that combines a pallet jack and a weighing scale. It allows you to weigh and move pallets at the same time, saving time and labor. Pallet jack scales are used in various industries, such as manufacturing, logistics, and shipping. Crane Scale : A crane scale is a device that measures the weight or mass of objects that are suspended from a crane. It has a hook at the bottom and a large display that allows distant viewing. Crane scales are used for various industrial applications, such as manufacturing, shipping, mining, recycling, and more Wheel Loader Scale : A wheel loader scale is a system that measures the weight of the materials lifted by a wheel loader, a type of heavy machinery used for moving large amounts of earth, sand, gravel, or other materials. A wheel loader scale can help improve the efficiency and accuracy of loading operations, as well as the inventory management and safety of the industries that use them. A wheel loader scale typically consists of a hydraulic sensor, a display unit, and a data management system. The hydraulic sensor is installed in the wheel loader and detects the pressure changes caused by the load. The display unit shows the weight information to the operator and allows them to set target loads, select products and customers, and export data. The data management system can store, analyze, and transmit the weight data to other devices or platforms. Testing and certification Most countries regulate the design and servicing of scales used for commerce. For example, in the European Union weighing instruments are subject to 2014/31/EU and 2014/32/EU directives. A conformity assessment procedure is carried out before placing the instrument on the market, andv the instruments are verified after a given period of time in member states of the European Union. This has tended to cause scale technology to lag behind other technologies because expensive regulatory hurdles are involved in introducing new designs. Nevertheless, there has been a trend to "digital load cells" which are actually strain-gauge cells with dedicated analog converters and networking built into the cell itself. Such designs have reduced the service problems inherent with combining and transmitting a number of 20 millivolt signals in hostile environments. Government regulation generally requires periodic inspections by licensed technicians, using masses whose calibration is traceable to an approved laboratory. Scales intended for non-trade use, such as those used in bathrooms, doctor's offices, kitchens (portion control), and price estimation (but not official price determination) may be produced, but must by law be labelled "Not Legal for Trade" to ensure that they are not re-purposed in a way that jeopardizes commercial interest. In the United States, the document describing how scales must be designed, installed, and used for commercial purposes is NIST Handbook 44. Legal For Trade (LFT) certification usually approve the readability by testing repeatability of measurements to ensure a maximum margin of error of 10%. Because gravity varies by over 0.5% over the surface of the earth, the distinction between force due to gravity and mass is relevant for accurate calibration of scales for commercial purposes. Usually, the goal is to measure the mass of the sample rather than its force due to gravity at that particular location. Traditional mechanical balance-beam scales intrinsically measured mass. But ordinary electronic scales intrinsically measure the gravitational force between the sample and the earth, i.e. the weight of the sample, which varies with location. So such a scale has to be re-calibrated after installation, for that specific location, in order to obtain an accurate indication of mass. Sources of error Some of the sources of error in weighing are: Buoyancy – Objects in air develop a buoyancy force that is directly proportional to the volume of air displaced. The difference in density of air due to barometric pressure and temperature creates errors. Error in the mass of reference weight Air gusts, even small ones, which push the scale up or down Friction in the moving components that causes the scale to reach equilibrium at a different configuration than a frictionless equilibrium should occur. Settling airborne dust contributing to the weight Mis-calibration over time, due to drift in the circuit's accuracy, or temperature change Mis-aligned mechanical components due to thermal expansion or contraction of components Magnetic fields acting on ferrous components Forces from electrostatic fields, for example, from feet shuffled on carpets on a dry day Chemical reactivity between air and the substance being weighed (or the balance itself, in the form of corrosion) Condensation of atmospheric water on cold items Evaporation of water from wet items Convection of air from hot or cold items Gravitational differences for a scale which measures force, but not for a balance. Vibration and seismic disturbances Hybrid spring and balance scales Elastic arm scale In 2014 a concept of hybrid scale was introduced, the elastically deformable arm scale, which is a combination between a spring scale and a beam balance, exploiting simultaneously both principles of equilibrium and deformation. In this scale, the rigid arms of a classical beam balance (for example a steelyard) are replaced with a flexible elastic rod in an inclined frictionless sliding sleeve. The rod can reach a unique sliding equilibrium when two vertical dead loads (or masses) are applied at its edges. Equilibrium, which would be impossible with rigid arms, is guaranteed because configurational forces develop at the two edges of the sleeve as a consequence of both the free sliding condition and the nonlinear kinematics of the elastic rod. This mass measuring device can also work without a counterweight.	Physical sciences	Classical mechanics	null
964475	https://en.wikipedia.org/wiki/Dwarf%20galaxy	Dwarf galaxy	A dwarf galaxy is a small galaxy composed of about 1000 up to several billion stars, as compared to the Milky Way's 200–400 billion stars. The Large Magellanic Cloud, which closely orbits the Milky Way and contains over 30 billion stars, is sometimes classified as a dwarf galaxy; others consider it a full-fledged galaxy. Dwarf galaxies' formation and activity are thought to be heavily influenced by interactions with larger galaxies. Astronomers identify numerous types of dwarf galaxies, based on their shape and composition. Formation One theory states that most galaxies, including dwarf galaxies, form in association with dark matter, or from gas that contains metals. However, NASA's Galaxy Evolution Explorer space probe identified new dwarf galaxies forming out of gases with low metallicity. These galaxies were located in the Leo Ring, a cloud of hydrogen and helium around two massive galaxies in the constellation Leo. Because of their small size, dwarf galaxies have been observed being pulled toward and ripped by neighbouring spiral galaxies, resulting in stellar streams and eventually galaxy merger. Local dwarf galaxies There are many dwarf galaxies in the Local Group; these small galaxies frequently orbit larger galaxies, such as the Milky Way, the Andromeda Galaxy and the Triangulum Galaxy. A 2007 paper has suggested that many dwarf galaxies were created by galactic tides during the early evolutions of the Milky Way and Andromeda. Tidal dwarf galaxies are produced when galaxies collide and their gravitational masses interact. Streams of galactic material are pulled away from the parent galaxies and the halos of dark matter that surround them. A 2018 study suggests that some local dwarf galaxies formed extremely early, during the Dark Ages within the first billion years after the Big Bang. More than 20 known dwarf galaxies orbit the Milky Way, and recent observations have also led astronomers to believe the largest globular cluster in the Milky Way, Omega Centauri, is in fact the core of a dwarf galaxy with a black hole at its centre, which was at some time absorbed by the Milky Way. Common types Elliptical galaxy: dwarf elliptical galaxy (dE) Dwarf spheroidal galaxy (dSph): Once a subtype of dwarf ellipticals, now regarded as a distinct type Irregular galaxy: dwarf irregular galaxy (dIrr) Spiral galaxy: dwarf spiral galaxy (dS) Magellanic type dwarfs Blue compact dwarf galaxies (see section below) Ultra-compact dwarf galaxies (see section below) Blue compact dwarf galaxies In astronomy, a blue compact dwarf galaxy (BCD galaxy) is a small galaxy which contains large clusters of young, hot, massive stars. These stars, the brightest of which are blue, cause the galaxy itself to appear blue in colour. Most BCD galaxies are also classified as dwarf irregular galaxies or as dwarf lenticular galaxies. Because they are composed of star clusters, BCD galaxies lack a uniform shape. They consume gas intensely, which causes their stars to become very violent when forming. BCD galaxies cool in the process of forming new stars. The galaxies' stars are all formed at different time periods, so the galaxies have time to cool and to build up matter to form new stars. As time passes, this star formation changes the shape of the galaxies. Nearby examples include NGC 1705, NGC 2915, NGC 3353 and UGCA 281. Ultra-faint dwarf galaxies Ultra-faint dwarf galaxies (UFDs) are a class of galaxies that contain from a few hundred to one hundred thousand stars, making them the faintest galaxies in the Universe. UFDs resemble globular clusters (GCs) in appearance but have very different properties. Unlike GCs, UFDs contain a significant amount of dark matter and are more extended. UFDs were first discovered with the advent of digital sky surveys in 2005, in particular with the Sloan Digital Sky Survey (SDSS). UFDs are the most dark matter-dominated systems known. Astronomers believe that UFDs encode valuable information about the early Universe, as all UFDs discovered so far are ancient systems that have likely formed very early on, only a few million years after the Big Bang and before the epoch of reionization. Recent theoretical work has hypothesised the existence of a population of young UFDs that form at a much later time than the ancient UFDs. These galaxies have not been observed in our Universe so far. Ultra-compact dwarfs Ultra-compact dwarf galaxies (UCD) are a class of very compact galaxies with very high stellar densities, discovered in the 2000s. They are thought to be on the order of 200 light years across, containing about 100 million stars. It is theorised that these are the cores of nucleated dwarf elliptical galaxies that have been stripped of gas and outlying stars by tidal interactions, travelling through the hearts of rich clusters. UCDs have been found in the Virgo Cluster, Fornax Cluster, Abell 1689, and the Coma Cluster, amongst others. In particular, an unprecedentedly large sample of ~ 100 UCDs has been found in the core region of the Virgo cluster by the Next Generation Virgo Cluster Survey team. The first ever relatively robust studies of the global properties of Virgo UCDs suggest that UCDs have distinct dynamical and structural properties from normal globular clusters. An extreme example of UCD is M60-UCD1, about 54 million light years away, which contains approximately 200 million solar masses within a 160 light year radius; the stars in its central region are packed 25 times more densely than stars in Earth's region in the Milky Way. M59-UCD3 is approximately the same size as M60-UCD1 with a half-light radius, rh, of approximately 20 parsecs but is 40% more luminous with an absolute visual magnitude of approximately −14.6. This makes M59-UCD3 the second densest known galaxy. Based on stellar orbital velocities, two UCD in the Virgo Cluster are claimed to have supermassive black holes weighing 13% and 18% of the galaxies' masses. Partial list Aquarius Dwarf Canis Major Dwarf Galaxy Carina Dwarf Crater 2 dwarf Draco Dwarf Eridanus II Fornax Dwarf Henize 2-10 I Zwicky 18 IC 10 Large Magellanic Cloud Leo I Leo II NGC 1569 NGC 1705 NGC 2915 NGC 3353 Pegasus Dwarf Irregular Galaxy PHL 293B Phoenix Dwarf Sagittarius Dwarf Spheroidal Galaxy Sagittarius Dwarf Irregular Galaxy Sculptor Dwarf Galaxy Sculptor Dwarf Irregular Galaxy Sextans A Sextans Dwarf Spheroidal Small Magellanic Cloud Tucana Dwarf Ursa Major I Dwarf Ursa Major II Dwarf Ursa Minor Dwarf Willman 1 Gallery	Physical sciences	Galaxy morphological classification	null
965247	https://en.wikipedia.org/wiki/Spanish%20dancer	Spanish dancer	The Spanish dancer, scientific name Hexabranchus sanguineus (literally meaning "blood-colored six-gills"), is a dorid nudibranch, a very large and colorful sea slug, a marine gastropod mollusk in the family Hexabranchidae.The taxonomy of the genus Hexabranchus has been controversial but a thorough molecular and morphological study published in 2023 showed that the name H. sanguineus was being used for at least 5 distinct species. Description Hexabranchus sanguineus is a large dorid nudibranch which commonly grows up to a maximum length of 25 cm, with some reports to 40 cm in the Red Sea. All Hexabranchus species have soft, flattened bodies, the anterior dorsal portion has a pair of retractable rhinophores and the posterior part has six contractile gills inserted independently in the body. The pair of oral tentacles are constituted by a fine flexible membrane provided with large digital lobes. In a normal situation when the animal is crawling, the edges of its mantle are curled upwards creating a peripheral blister. If the animal is disturbed, it unfolds its edges and can swim through contractions and undulations of the body to move away from the disturbing element. Its common name, Spanish dancer, comes from this particular defense. Distribution and habitat This species is reported from the Red Sea, the Western Indian Ocean, French Polynesia and the Western Pacific, with different colour morphs in each region which are not differentiated by morphology or DNA barcodes. It likes rocky and coral reefs with many sponges and shelters from 1 to 50 meters deep. Biology During daytime, the Spanish dancer hides away from the light in the crevices of its natural habitat to only come out late at night. It feeds on various species of sponge. Like all nudibranchs, it is hermaphrodite and its bright red to pink egg ribbon has a spiral shape related to the size of the animal so relatively large. The Spanish Dancer consumes sponges from the family Halichondriidae. Once consumed, the Spanish Dancer derives a potent chemical that it can use as defense. Hexabranchus sanguineus then passes the defensive compounds obtained into its egg ribbons via macrolides, giving the physically defenseless egg ribbons a toxin defense. The latter is coveted by some other species of nudibranch as Favorinus tsuruganus or Favorinus japonicus. The emperor shrimp, Periclimenes imperator, is a commensal shrimp that is commonly found living on Hexabranchus sanguineus.	Biology and health sciences	Gastropods	Animals
965323	https://en.wikipedia.org/wiki/Antimicrobial	Antimicrobial	An antimicrobial is an agent that kills microorganisms (microbicide) or stops their growth (bacteriostatic agent). Antimicrobial medicines can be grouped according to the microorganisms they act primarily against. For example, antibiotics are used against bacteria, and antifungals are used against fungi. They can also be classified according to their function. Antimicrobial medicines to treat infection are known as ⠀⠀antimicrobial chemotherapy, while antimicrobial drugs are used to prevent infection, which known as antimicrobial prophylaxis. The main classes of antimicrobial agents are disinfectants (non-selective agents, such as bleach), which kill a wide range of microbes on non-living surfaces to prevent the spread of illness, antiseptics (which are applied to living tissue and help reduce infection during surgery), and antibiotics (which destroy microorganisms within the body). The term antibiotic originally described only those formulations derived from living microorganisms but is now also applied to synthetic agents, such as sulfonamide's or fluoroquinolone's. Though the term used to be restricted to antibacterial, and is often used as a synonym for them by medical professionals and in medical literature, its context has broadened to include all antimicrobials. Antibacterial agents can be further subdivided into bactericidal agents, which kill bacteria, and bacteriostatic agents, which slow down or stall bacterial growth. In response, further advancements in antimicrobial technologies have resulted in solutions that can go beyond simply inhibiting microbial growth. Instead, certain types of porous media have been developed to kill microbes on contact. The misuse and overuse of antimicrobials in humans, animals and plants are the main drivers in the development of drug-resistant pathogens. It is estimated that bacterial antimicrobial resistance (AMR) was directly responsible for 1.27 million global deaths in 2019 and contributed to 4.95 million deaths. History Antimicrobial use has been common practice for at least 2000 years. Ancient Egyptians and ancient Greeks used specific molds and plant extracts to treat infection. In the 19th century, microbiologists such as Louis Pasteur and Jules Francois Joubert observed antagonism between some bacteria and discussed the merits of controlling these interactions in medicine. Louis Pasteur's work in fermentation and spontaneous generation led to the distinction between anaerobic and aerobic bacteria. The information garnered by Pasteur led Joseph Lister to incorporate antiseptic methods, such as sterilizing surgical tools and debriding wounds into surgical procedures. The implementation of these antiseptic techniques drastically reduced the number of infections and subsequent deaths associated with surgical procedures. Louis Pasteur's work in microbiology also led to the development of many vaccines for life-threatening diseases such as anthrax and rabies. On September 3, 1928, Alexander Fleming returned from a vacation and discovered that a Petri dish filled with Staphylococcus was separated into colonies due to the antimicrobial fungus Penicillium rubens. Fleming and his associates struggled to isolate the antimicrobial but referenced its therapeutic potential in 1929 in the British Journal of Experimental Pathology. In 1942, Howard Florey, Ernst Chain, and Edward Abraham used Fleming's work to purify and extract penicillin for medicinal uses earning them the 1945 Nobel Prize in Medicine. Chemical Antibacterials Antibacterials are used to treat bacterial infections. Antibiotics are classified generally as beta-lactams, macrolides, quinolones, tetracyclines or aminoglycosides. Their classification within these categories depends on their antimicrobial spectra, pharmacodynamics and chemical composition. Prolonged use of certain antibacterials can decrease the number of enteric bacteria, which may have a negative impact on health. Consumption of probiotics and healthy eating may help to replace destroyed gut flora. Stool transplants may be considered however for patients who are having difficulty recovering from prolonged antibiotic treatment, such as recurrent Clostridioides difficile infections. The discovery, development and use of antibacterials during the 20th century have reduced mortality from bacterial infections. The antibiotic era began with the therapeutic application of sulfonamide drugs in 1936, followed by a "golden" period of discovery from about 1945 to 1970, when a number of structurally diverse and highly effective agents were discovered and developed. Since 1980, the introduction of new antimicrobial agents for clinical use has declined, in part because of the enormous expense of developing and testing new drugs. In parallel, there has been an alarming increase in antimicrobial resistance of bacteria, fungi, parasites and some viruses to multiple existing agents. Antibacterials are among the most commonly used and misused drugs by physicians, for example, in viral respiratory tract infections. As a consequence of widespread and injudicious use of antibacterials, there has been an accelerated emergence of antibiotic-resistant pathogens, resulting in a serious threat to global public health. The resistance problem demands that a renewed effort be made to seek antibacterial agents effective against pathogenic bacteria resistant to current antibacterials. Possible strategies towards this objective include increased sampling from diverse environments and application of metagenomics to identify bioactive compounds produced by currently unknown and uncultured microorganisms as well as the development of small-molecule libraries customized for bacterial targets. Antifungals Antifungals are used to kill or prevent further growth of fungi. In medicine, they are used as a treatment for infections such as athlete's foot, ringworm and thrush and work by exploiting differences between mammalian and fungal cells. Unlike bacteria, both fungi and humans are eukaryotes. Thus, fungal and human cells are similar at the molecular level, making it more difficult to find a target for an antifungal drug to attack that does not also exist in the host organism. Consequently, there are often side effects to some of these drugs. Some of these side effects can be life-threatening if the drug is not used properly. As well as their use in medicine, antifungals are frequently sought after to control indoor mold in damp or wet home materials. Sodium bicarbonate (baking soda) blasted on to surfaces acts as an antifungal. Another antifungal solution applied after or without blasting by soda is a mix of hydrogen peroxide and a thin surface coating that neutralizes mold and encapsulates the surface to prevent spore release. Some paints are also manufactured with an added antifungal agent for use in high humidity areas such as bathrooms or kitchens. Other antifungal surface treatments typically contain variants of metals known to suppress mold growth e.g. pigments or solutions containing copper, silver or zinc. These solutions are not usually available to the general public because of their toxicity. Antivirals Antiviral drugs are a class of medication used specifically for treating viral infections. Like antibiotics, specific antivirals are used for specific viruses. They should be distinguished from viricides, which actively deactivate virus particles outside the body. Many antiviral drugs are designed to treat infections by retroviruses, including HIV. Important antiretroviral drugs include the class of protease inhibitors. Herpes viruses, best known for causing cold sores and genital herpes, are usually treated with the nucleoside analogue acyclovir. Viral hepatitis is caused by five unrelated hepatotropic viruses (A-E) and may be treated with antiviral drugs depending on the type of infection. Some influenza A and B viruses have become resistant to neuraminidase inhibitors such as oseltamivir, and the search for new substances continues. Antiparasitics Antiparasitics are a class of medications indicated for the treatment of infectious diseases such as leishmaniasis, malaria and Chagas disease, which are caused by parasites such as nematodes, cestodes, trematodes and infectious protozoa. Antiparasitic medications include metronidazole, iodoquinol and albendazole. Like all therapeutic antimicrobials, they must kill the infecting organism without serious damage to the host. Broad-spectrum therapeutics Broad-spectrum therapeutics are active against multiple classes of pathogens. Such therapeutics have been suggested as potential emergency treatments for pandemics. Non-pharmaceutical A wide range of chemical and natural compounds are used as antimicrobials. Organic acids and their salts are used widely in food products, e.g. lactic acid, citric acid, acetic acid, either as ingredients or as disinfectants. For example, beef carcasses often are sprayed with acids, and then rinsed or steamed, to reduce the prevalence of Escherichia coli. Heavy metal cations such as Hg2+ and Pb2+ have antimicrobial activities, but can be toxic. In recent years, the antimicrobial activity of coordination compounds has been investigated. Traditional herbalists used plants to treat infectious disease. Many of these plants have been investigated scientifically for antimicrobial activity, and some plant products have been shown to inhibit the growth of pathogenic microorganisms. A number of these agents appear to have structures and modes of action that are distinct from those of the antibiotics in current use, suggesting that cross-resistance with agents already in use may be minimal. Copper Copper-alloy surfaces have natural intrinsic antimicrobial properties and can kill microorganisms such as E. coli and Staphylococcus. The United States Environmental Protection Agency approved the registration of antimicrobial copper alloy surfaces for use in addition to regular cleaning and disinfection to control infections. Antimicrobial copper alloys are being installed in some healthcare facilities and subway transit systems as a public hygienic measure. Copper nanoparticles are attracting interest for the intrinsic antimicrobial behaviours. Essential oils Many essential oils included in herbal pharmacopoeias are claimed to possess antimicrobial activity, with the oils of bay, cinnamon, clove and thyme reported to be the most potent in studies with foodborne bacterial pathogens. Coconut oil is also known for its antimicrobial properties. Active constituents include terpenoids and secondary metabolites. Despite their prevalent use in alternative medicine, essential oils have seen limited use in mainstream medicine. While 25 to 50% of pharmaceutical compounds are plant-derived, none are used as antimicrobials, though there has been increased research in this direction. Barriers to increased usage in mainstream medicine include poor regulatory oversight and quality control, mislabeled or misidentified products, and limited modes of delivery. Antimicrobial pesticides According to the U.S. Environmental Protection Agency (EPA), and defined by the Federal Insecticide, Fungicide, and Rodenticide Act, antimicrobial pesticides are used to control growth of microbes through disinfection, sanitation, or reduction of development and to protect inanimate objects, industrial processes or systems, surfaces, water, or other chemical substances from contamination, fouling, or deterioration caused by bacteria, viruses, fungi, protozoa, algae, or slime. The EPA monitors products, such as disinfectants/sanitizers for use in hospitals or homes, to ascertain efficacy. Products that are meant for public health are therefore under this monitoring system, including products used for drinking water, swimming pools, food sanitation, and other environmental surfaces. These pesticide products are registered under the premise that, when used properly, they do not demonstrate unreasonable side effects to humans or the environment. Even once certain products are on the market, the EPA continues to monitor and evaluate them to make sure they maintain efficacy in protecting public health. Public health products regulated by the EPA are divided into three categories: Disinfectants: Destroy or inactivate microorganisms (bacteria, fungi, viruses,) but may not act as sporicides (as those are the most difficult form to destroy). According to efficacy data, the EPA will classify a disinfectant as limited, general/ broad spectrum, or as a hospital disinfectant. Sanitizers: Reduce the number of microorganisms, but may not kill or eliminate all of them. Sterilizers (Sporicides): Eliminate all bacteria, fungi, spores, and viruses. Antimicrobial pesticide safety Antimicrobial pesticides have the potential to be a major factor in drug resistance. Organizations such as the World Health Organization call for significant reduction in their use globally to combat this. According to a 2010 Centers for Disease Control and Prevention report, health-care workers can take steps to improve their safety measures against antimicrobial pesticide exposure. Workers are advised to minimize exposure to these agents by wearing personal protective equipment such as gloves and safety glasses. Additionally, it is important to follow the handling instructions properly, as that is how the EPA has deemed them as safe to use. Employees should be educated about the health hazards and encouraged to seek medical care if exposure occurs. Ozone Ozone can kill microorganisms in air, water and process equipment and has been used in settings such as kitchen exhaust ventilation, garbage rooms, grease traps, biogas plants, wastewater treatment plants, textile production, breweries, dairies, food and hygiene production, pharmaceutical industries, bottling plants, zoos, municipal drinking-water systems, swimming pools and spas, and in the laundering of clothes and treatment of in–house mold and odors. Antimicrobial scrubs Antimicrobial scrubs can reduce the accumulation of odors and stains on scrubs, which in turn improves their longevity. These scrubs also come in a variety of colors and styles. As antimicrobial technology develops at a rapid pace, these scrubs are readily available, with more advanced versions hitting the market every year. These bacteria could then be spread to office desks, break rooms, computers, and other shared technology. This can lead to outbreaks and infections like methicillin-resistant staphylococcus aureus, treatments for which cost the healthcare industry $20 billion a year. Halogens Elements such as chlorine, iodine, fluorine, and bromine are nonmetallic in nature and constitute the halogen family. Each of these halogens have a different antimicrobial effect that is influenced by various factors such as pH, temperature, contact time, and type of microorganism. Chlorine and iodine are the two most commonly used antimicrobials. Chlorine is extensively used as a disinfectant in the water treatment plants, drug, and food industries. In wastewater treatment plants, chlorine is widely used as a disinfectant. It oxidizes soluble contaminants and kills bacteria and viruses. It is also highly effective against bacterial spores. The mode of action is by breaking the bonds present in these microorganisms. When a bacterial enzyme comes in contact with a compound containing chlorine, the hydrogen atom in that molecule gets displaced and is replaced with chlorine. This in turn changes the enzyme function which ultimately leads to the death of the bacterium. Iodine is most commonly used for sterilization and wound cleaning. The three major antimicrobial compounds containing iodine are alcohol-iodine solution, an aqueous solution of iodine, and iodophors. Iodophors are more bactericidal and are used as antiseptics as they are less irritating when applied to the skin. Bacterial spores on the other hand cannot be killed by iodine, but they can be inhibited by iodophors. The growth of microorganisms is inhibited when iodine penetrates into the cells and oxidizes proteins, genetic material, and fatty acids. Bromine is also an effective antimicrobial that is used in water treatment plants. When mixed with chlorine it is highly effective against bacterial spores such as S. faecalis. Alcohols Alcohols are commonly used as disinfectants and antiseptics. Alcohols kill vegetative bacteria, most viruses and fungi. Ethyl alcohol, n-propanol and isopropyl alcohol are the most commonly used antimicrobial agents. Methanol is also a disinfecting agent but is not generally used as it is highly poisonous. Escherichia coli, Salmonella, and Staphylococcus aureus are a few bacteria whose growth can be inhibited by alcohols. Alcohols have a high efficiency against enveloped viruses (60–70% ethyl alcohol) 70% isopropyl alcohol or ethanol are highly effective as an antimicrobial agent. In the presence of water, 70% alcohol causes coagulation of the proteins thus inhibiting microbial growth. Alcohols are not quite efficient when it comes to spores. The mode of action is by denaturing the proteins. Alcohols interfere with the hydrogen bonds present in the protein structure. Alcohols also dissolve the lipid membranes that are present in microorganisms. Disruption of the cell membrane is another property of alcohols that aids in cell death. Alcohols are cheap and effective antimicrobials. They are widely used in the pharmaceutical industry. Alcohols are commonly used in hand sanitizers, antiseptics, and disinfectants. Phenol and Phenolic compounds Phenol, also known as carbolic acid, was one of the first chemicals which was used as an antimicrobial agent. It has high antiseptic properties. It is bacteriostatic at concentrations of 0.1%–1% and is bactericidal/fungicidal at 1%–2%. A 5% solution kills anthrax spores in 48 hr. Phenols are most commonly used in oral mouth washes and household cleaning agents. They are active against a wide range of bacteria, fungi and viruses. Today phenol derivatives such as thymol and cresol are used because they are less toxic compared to phenol. These phenolic compounds have a benzene ring along with the –OH group incorporated into their structures. They have a higher antimicrobial activity. These compounds inhibit microbial growth by precipitating proteins which lead to their denaturation and by penetrating into the cell membrane of microorganisms and disrupting it. Phenolic compounds can also deactivate enzymes and damage the amino acids in microbial cells. Phenolics such as fentichlore, an antibacterial and antifungal agent, are used as an oral treatment for fungal infections. Trischlosan is highly effective against both gram-positive and gram-negative bacteria. Hexachlorophene (Bisphenol) is used as a surfactant. It is widely used in soaps, handwashes, and skin products because of its antiseptic properties. It is also used as a sterilizing agent. Cresol is an effective antimicrobial and is widely used in mouthwashes and cough drops. Phenolics have high antimicrobial activity against bacteria such as Staphylococcus epidermidis and Pseudomonas aeruginosa. 2-Phenylphenol-water solutions are used in immersion treatments of fruit for packing. (It is not used on the packing materials however.) Ihloff and Kalitzki 1961 find a small but measurable amount remains in the skin of fruits processed in this manner. Aldehydes Aldehydes are highly effective against bacteria, fungi, and viruses. Aldehydes inhibit bacterial growth by disrupting the outer membrane. They are used in the disinfection and sterilization of surgical instruments. As they are highly toxic, they are not used in antiseptics. Currently, only three aldehyde compounds are of widespread practical use as disinfectant biocides, namely glutaraldehyde, formaldehyde, and ortho-phthalaldehyde (OPA) despite the demonstration that many other aldehydes possess good antimicrobial activity. However, due to its long contact time other disinfectants are commonly preferred. Physical Heat Microorganisms have a minimum temperature, an optimum, and a maximum temperature for growth. High temperature as well as low temperatures are used as physical agents of control. Different organisms show different degrees of resistance or susceptibility to heat or temperature, some organisms such as bacterial endospore are more resistant while vegetative cells are less resistant and are easily killed at lower temperatures. Another method that involves the use of heat to kill microorganisms is fractional sterilization. This process involves the exposure to a temperature of 100 degrees Celsius for an hour per day for several days. Fractional sterilization is also called tyndallization. Bacterial endospores can be killed using this method. Both dry and moist heat are effective in eliminating microbial life. For example, jars used to store preserves such as jam can be sterilized by heating them in a conventional oven. Heat is also used in pasteurization, a method for slowing the spoilage of foods such as milk, cheese, juices, wines and vinegar. Such products are heated to a certain temperature for a set period of time, which greatly reduces the number of harmful microorganisms. Low temperature is also used to inhibit microbial activity by slowing down microbial metabolism. Radiation Foods are often irradiated to kill harmful pathogens. There are two types of radiations that are used to inhibit the growth of microorganisms – ionizing and non-ionizing radiations. Common sources of radiation used in food sterilization include cobalt-60 (a gamma emitter), electron beams and . Ultraviolet light is also used to disinfect drinking water, both in small-scale personal-use systems and larger-scale community water purification systems. Desiccation Desiccation is also known as dehydration. It is the state of extreme dryness or the process of extreme drying. Some microorganisms like bacteria, yeasts and molds require water for their growth. Desiccation dries up the water content thus inhibiting microbial growth. On the availability of water, the bacteria resume their growth, thus desiccation does not completely inhibit bacterial growth. The instrument used to carry out this process is called a desiccator. This process is widely used in the food industry and is an efficient method for food preservation. Desiccation is also largely used in the pharmaceutical industry to store vaccines and other products. Antimicrobial surfaces Antimicrobial surfaces are designed to either inhibit the ability of microorganisms to grow or damaging them by chemical (copper toxicity) or physical processes (micro/nano-pillars to rupture cell walls). These surfaces are especially important for the healthcare industry. Designing effective antimicrobial surfaces demands an in-depth understanding of the initial microbe-surface adhesion mechanisms. Molecular dynamics simulation and time-lapse imaging are typically used to investigate these mechanisms. Osmotic pressure Osmotic pressure is the pressure required to prevent a solvent from passing from a region of high concentration to a region of low concentration through a semipermeable membrane. When the concentration of dissolved materials or solute is higher inside the cell than it is outside, the cell is said to be in a hypotonic environment and water will flow into the cell.When the bacteria is placed in hypertonic solution, it causes plasmolysis or cell shrinking, similarly in hypotonic solution, bacteria undergoes plasmotysis or turgid state. This plasmolysis and plasmotysis kills bacteria because it causes change in osmotic pressure. Antimicrobial resistance Antimicrobial resistance The misuse and overuse of antimicrobials in humans, animals and plants are the main drivers in the development of drug-resistant pathogens. It is estimated that bacterial antimicrobial resistance (AMR) was directly responsible for 1.27 million global deaths in 2019 and contributed to 4.95 million deaths.	Biology and health sciences	Anti-infectives	Health
965924	https://en.wikipedia.org/wiki/Roller%20%28agricultural%20tool%29	Roller (agricultural tool)	The roller is an agricultural tool used for flattening land or breaking up large clumps of soil, especially after ploughing or disc harrowing. Typically, rollers are pulled by tractors or, prior to mechanisation, a team of animals such as horses or oxen. As well as for agricultural purposes, rollers are used on cricket pitches and residential lawn areas. Flatter land makes subsequent weed control and harvesting easier, and rolling can help to reduce moisture loss from cultivated soil. On lawns, rolling levels the land for mowing and compacts the soil surface. Rollers may be weighted in different ways. For many uses a heavy roller is used. These may consist of one or more cylinders made of thick steel, a thinner steel cylinder filled with concrete, or a cylinder filled with water. A water-filled roller has the advantage that the water may be drained out for lighter use or for transport. In frost-prone areas a water filled roller must be drained for winter storage to avoid breakage due to the expansion for water as it turns to ice. Designs One-piece versus segmented On tilled soil a one-piece roller has the disadvantage that when turning corners the outer end of the roller has to rotate much faster than the inner end, forcing one or both ends to skid. A one-piece roller turned on soft ground will skid up a heap of soil at the outer radius, leaving heaps, which is counter-productive. Rollers are often made in two or three sections to reduce this problem, and the Cambridge roller overcomes it altogether by mounting many small segments onto one axle so that they can each rotate at local ground-speed. Smooth versus ridged The surface of rollers may be smooth, or it may be textured to help break up soil or to groove the final surface to reduce scouring from rain. Each segment of a Cambridge roller has a rib around its edge for this purpose. The name cultipacker is often used for such ridged types, especially in the United States. Uses Farming use Rollers are a secondary tillage tool used for flattening land or breaking up large clumps of soil, especially after ploughing or disc harrowing. Rollers are typically pulled by tractors today. Before mechanised agriculture, a team of working animals such as horses or oxen provided the power. Animal power is still used today in some contexts, such as on Amish farms in the United States and in regions of Asia where draft oxen are still widely used. Rollers prepare optimal seedbeds by making them as flat as is practical and moderately firmed. Flatness is important at planting because it is the only practical way to control average seed planting depth without laborious hand planting of each seed; it is not practical to follow an instruction of (for example) 1-cm planting depth if the contour of the seedbed varies by 2 cm or more between adjacent spots. This is why breaking up of even small clods/lumps, and well-leveled spreading of soil, is important at planting time. Flatter land also makes subsequent weed control and harvesting easier. For example, in mechanical weed control, controlling cultivator tooth depth is practical only with a decently flat soil contour, and in combining, controlling combine head height is practical only with a decently flat soil contour. Rolling is also believed to help reduce moisture loss from cultivated soil. Ganging and trailing Rollers may be ganged to increase the width of each pass/swath. Rollers may be trailed after other equipment such as ploughs, disc harrows, or mowers. Cricket pitch use In cricket, rollers are used to make the pitch flat and less dangerous for batsmen. Several size rollers have been used in the history of cricket, from light rollers that were used in the days of uncovered pitches and at some stages during the 1950s to make batting less easy, to the modern “heavy roller” universally used in top-class cricket today. Regulations permit a pitch only to be rolled at the commencement of each innings or day’s play, but this has still had a massive influence on the game by eliminating the shooters that were ubiquitous on all but light soils before heavy rollers were used. Heavy rollers have sometimes been criticised for making batting too easy and for reducing the rate at which pitches dry out after rain in the cool English climate. Lawn use Lawn rollers are designed to even out or firm up the lawn surface, especially in climates where heaving causes the lawn to be lumpy. Heaving may result when the ground freezes and thaws many times over winter. Where this occurs, gardeners are advised to give the lawn a light rolling with a lawn roller in the spring. Clay or wet soils should not be rolled as they become compacted.	Technology	Agricultural tools	null
11694119	https://en.wikipedia.org/wiki/Inverse%20second	Inverse second	The inverse second or reciprocal second (s−1), also called per second, is a unit defined as the multiplicative inverse of the second (a unit of time). It is applicable for physical quantities of dimension reciprocal time, such as frequency and strain rate. It is dimensionally equivalent to: hertz (Hz), historically known as cycles per second – the SI unit for frequency and rotational frequency becquerel (Bq) – the SI unit for the rate of occurrence of aperiodic or stochastic radionuclide events baud (Bd) – the unit for symbol rate over a communication link bit per second (bit/s) – the unit of bit rate However, the special names and symbols above for s−1 are recommend for clarity. Reciprocal second should not be confused with radian per second (rad⋅s−1), the SI unit for angular frequency and angular velocity. As the radian is a dimensionless unit, radian per second is dimensionally consistent with reciprocal second. However, they are used for different kinds of quantity, frequency and angular frequency, whose numerical value differs by 2. The inverse minute or reciprocal minute (min−1), also called per minute, is 60−1 s−1, as 1 min = 60 s; it is used in quantities of type "counts per minute", such as: Actions per minute Beats per minute Counts per minute Revolutions per minute (rpm) Words per minute Inverse square second (s−2) is involved in the units of linear acceleration, angular acceleration, and rotational acceleration.	Physical sciences	Angular velocity	Basics and measurement
562526	https://en.wikipedia.org/wiki/Sugar%20alcohol	Sugar alcohol	Sugar alcohols (also called polyhydric alcohols, polyalcohols, alditols or glycitols) are organic compounds, typically derived from sugars, containing one hydroxyl group attached to each carbon atom. They are white, water-soluble solids that can occur naturally or be produced industrially by hydrogenating sugars. Since they contain multiple groups, they are classified as polyols. Sugar alcohols are used widely in the food industry as thickeners and sweeteners. In commercial foodstuffs, sugar alcohols are commonly used in place of table sugar (sucrose), often in combination with high-intensity artificial sweeteners, in order to offset their low sweetness. Xylitol and sorbitol are popular sugar alcohols in commercial foods. Structure Sugar alcohols have the general formula . In contrast, sugars have two fewer hydrogen atoms, for example, or . Like their parent sugars, sugar alcohols exist in diverse chain length. Most have five- or six-carbon chains, because they are derived respectively from pentoses (five-carbon sugars) and hexoses (six-carbon sugars), which are the more common sugars. They have one −OH group attached to each carbon. They are further differentiated by the relative orientation (stereochemistry) of these −OH groups. Unlike sugars, which tend to exist as rings, sugar alcohols do not, although they can be dehydrated to give cyclic ethers (e.g. sorbitan can be dehydrated to isosorbide). Production Sugar alcohols can be, and often are, produced from renewable resources. Particular feedstocks are starch, cellulose and hemicellulose; the main conversion technologies use as the reagent: hydrogenolysis, i.e. the cleavage of single bonds, converting polymers to smaller molecules, and hydrogenation of double bonds, converting sugars to sugar alcohols. Sorbitol and mannitol Mannitol is no longer obtained from natural sources; currently, sorbitol and mannitol are obtained by hydrogenation of sugars, using Raney nickel catalysts. The conversion of glucose and mannose to sorbitol and mannitol is given as Erythritol Erythritol is obtained by the fermentation of glucose and sucrose. Health effects Sugar alcohols do not contribute to tooth decay; in fact, xylitol deters tooth decay. Sugar alcohols are absorbed at 50% of the rate of sugars, resulting in less of an effect on blood sugar levels as measured by comparing their effect to sucrose using the glycemic index. Common sugar alcohols Ethylene glycol (2-carbon) Glycerol (3-carbon) Erythritol (4-carbon) Threitol (4-carbon) Arabitol (5-carbon) Xylitol (5-carbon) Ribitol (5-carbon) Mannitol (6-carbon) Sorbitol (6-carbon) Galactitol (6-carbon) Fucitol (6-carbon) Iditol (6-carbon) Inositol (6-carbon; a cyclic sugar alcohol) Volemitol (7-carbon) Isomalt (12-carbon) Maltitol (12-carbon) Lactitol (12-carbon) Maltotriitol (18-carbon) Maltotetraitol (24-carbon) Polyglycitol Both disaccharides and monosaccharides can form sugar alcohols; however, sugar alcohols derived from disaccharides (e.g. maltitol and lactitol) are not entirely hydrogenated because only one aldehyde group is available for reduction. Sugar alcohols as food additives This table presents the relative sweetness and food energy of the most widely used sugar alcohols. Despite the variance in food energy content of sugar alcohols, the European Union's labeling requirements assign a blanket value of 2.4 kcal/g to all sugar alcohols. Characteristics As a group, sugar alcohols are not as sweet as sucrose, and they have slightly less food energy than sucrose. Their flavor is similar to sucrose, and they can be used to mask the unpleasant aftertastes of some high-intensity sweeteners. Sugar alcohols are not metabolized by oral bacteria, and so they do not contribute to tooth decay. They do not brown or caramelize when heated. In addition to their sweetness, some sugar alcohols can produce a noticeable cooling sensation in the mouth when highly concentrated, for instance in sugar-free hard candy or chewing gum. This happens, for example, with the crystalline phase of sorbitol, erythritol, xylitol, mannitol, lactitol and maltitol. The cooling sensation is due to the dissolution of the sugar alcohol being an endothermic (heat-absorbing) reaction, one with a strong heat of solution. Absorption from the small intestine Sugar alcohols are usually incompletely absorbed into the blood stream from the small intestine which generally results in a smaller change in blood glucose than "regular" sugar (sucrose). This property makes them popular sweeteners among diabetics and people on low-carbohydrate diets. As an exception, erythritol is actually absorbed in the small intestine and excreted unchanged through urine, so it contributes no calories even though it is rather sweet. Side effects Like many other incompletely digestible substances, overconsumption of sugar alcohols can lead to bloating, diarrhea and flatulence because they are not fully absorbed in the small intestine. Some individuals experience such symptoms even in a single-serving quantity. With continued use, most people develop a degree of tolerance to sugar alcohols and no longer experience these symptoms.	Physical sciences	Sugar alcohols	Chemistry
562782	https://en.wikipedia.org/wiki/Vertex%20cover	Vertex cover	In graph theory, a vertex cover (sometimes node cover) of a graph is a set of vertices that includes at least one endpoint of every edge of the graph. In computer science, the problem of finding a minimum vertex cover is a classical optimization problem. It is NP-hard, so it cannot be solved by a polynomial-time algorithm if P ≠ NP. Moreover, it is hard to approximate – it cannot be approximated up to a factor smaller than 2 if the unique games conjecture is true. On the other hand, it has several simple 2-factor approximations. It is a typical example of an NP-hard optimization problem that has an approximation algorithm. Its decision version, the vertex cover problem, was one of Karp's 21 NP-complete problems and is therefore a classical NP-complete problem in computational complexity theory. Furthermore, the vertex cover problem is fixed-parameter tractable and a central problem in parameterized complexity theory. The minimum vertex cover problem can be formulated as a half-integral, linear program whose dual linear program is the maximum matching problem. Vertex cover problems have been generalized to hypergraphs, see Vertex cover in hypergraphs. Definition Formally, a vertex cover of an undirected graph is a subset of such that , that is to say it is a set of vertices where every edge has at least one endpoint in the vertex cover . Such a set is said to cover the edges of . The upper figure shows two examples of vertex covers, with some vertex cover marked in red. A minimum vertex cover is a vertex cover of smallest possible size. The vertex cover number is the size of a minimum vertex cover, i.e. . The lower figure shows examples of minimum vertex covers in the previous graphs. Examples The set of all vertices is a vertex cover. The endpoints of any maximal matching form a vertex cover. The complete bipartite graph has a minimum vertex cover of size . Properties A set of vertices is a vertex cover if and only if its complement is an independent set. Consequently, the number of vertices of a graph is equal to its minimum vertex cover number plus the size of a maximum independent set. Computational problem The minimum vertex cover problem is the optimization problem of finding a smallest vertex cover in a given graph. INSTANCE: Graph OUTPUT: Smallest number such that has a vertex cover of size . If the problem is stated as a decision problem, it is called the vertex cover problem: INSTANCE: Graph and positive integer . QUESTION: Does have a vertex cover of size at most ? The vertex cover problem is an NP-complete problem: it was one of Karp's 21 NP-complete problems. It is often used in computational complexity theory as a starting point for NP-hardness proofs. ILP formulation Assume that every vertex has an associated cost of . The (weighted) minimum vertex cover problem can be formulated as the following integer linear program (ILP). {\| \| minimize \| colspan="2" \| \| \| (minimize the total cost) \|- \| subject to \| \| for all \| \| (cover every edge of the graph), \|- \| \| \| for all . \| \| (every vertex is either in the vertex cover or not) \|} This ILP belongs to the more general class of ILPs for covering problems. The integrality gap of this ILP is , so its relaxation (allowing each variable to be in the interval from 0 to 1, rather than requiring the variables to be only 0 or 1) gives a factor- approximation algorithm for the minimum vertex cover problem. Furthermore, the linear programming relaxation of that ILP is half-integral, that is, there exists an optimal solution for which each entry is either 0, 1/2, or 1. A 2-approximate vertex cover can be obtained from this fractional solution by selecting the subset of vertices whose variables are nonzero. Exact evaluation The decision variant of the vertex cover problem is NP-complete, which means it is unlikely that there is an efficient algorithm to solve it exactly for arbitrary graphs. NP-completeness can be proven by reduction from 3-satisfiability or, as Karp did, by reduction from the clique problem. Vertex cover remains NP-complete even in cubic graphs and even in planar graphs of degree at most 3. For bipartite graphs, the equivalence between vertex cover and maximum matching described by Kőnig's theorem allows the bipartite vertex cover problem to be solved in polynomial time. For tree graphs, an algorithm finds a minimal vertex cover in polynomial time by finding the first leaf in the tree and adding its parent to the minimal vertex cover, then deleting the leaf and parent and all associated edges and continuing repeatedly until no edges remain in the tree. Fixed-parameter tractability An exhaustive search algorithm can solve the problem in time 2knO(1), where k is the size of the vertex cover. Vertex cover is therefore fixed-parameter tractable, and if we are only interested in small k, we can solve the problem in polynomial time. One algorithmic technique that works here is called bounded search tree algorithm, and its idea is to repeatedly choose some vertex and recursively branch, with two cases at each step: place either the current vertex or all its neighbours into the vertex cover. The algorithm for solving vertex cover that achieves the best asymptotic dependence on the parameter runs in time . The klam value of this time bound (an estimate for the largest parameter value that could be solved in a reasonable amount of time) is approximately 190. That is, unless additional algorithmic improvements can be found, this algorithm is suitable only for instances whose vertex cover number is 190 or less. Under reasonable complexity-theoretic assumptions, namely the exponential time hypothesis, the running time cannot be improved to 2o(k), even when is . However, for planar graphs, and more generally, for graphs excluding some fixed graph as a minor, a vertex cover of size k can be found in time , i.e., the problem is subexponential fixed-parameter tractable. This algorithm is again optimal, in the sense that, under the exponential time hypothesis, no algorithm can solve vertex cover on planar graphs in time . Approximate evaluation One can find a factor-2 approximation by repeatedly taking both endpoints of an edge into the vertex cover, then removing them from the graph. Put otherwise, we find a maximal matching M with a greedy algorithm and construct a vertex cover C that consists of all endpoints of the edges in M. In the following figure, a maximal matching M is marked with red, and the vertex cover C is marked with blue. The set C constructed this way is a vertex cover: suppose that an edge e is not covered by C; then M ∪ {e} is a matching and e ∉ M, which is a contradiction with the assumption that M is maximal. Furthermore, if e = {u, v} ∈ M, then any vertex cover – including an optimal vertex cover – must contain u or v (or both); otherwise the edge e is not covered. That is, an optimal cover contains at least one endpoint of each edge in M; in total, the set C is at most 2 times as large as the optimal vertex cover. This simple algorithm was discovered independently by Fanica Gavril and Mihalis Yannakakis. More involved techniques show that there are approximation algorithms with a slightly better approximation factor. For example, an approximation algorithm with an approximation factor of is known. The problem can be approximated with an approximation factor in - dense graphs. Inapproximability No better constant-factor approximation algorithm than the above one is known. The minimum vertex cover problem is APX-complete, that is, it cannot be approximated arbitrarily well unless P = NP. Using techniques from the PCP theorem, Dinur and Safra proved in 2005 that minimum vertex cover cannot be approximated within a factor of 1.3606 for any sufficiently large vertex degree unless P = NP. Later, the factor was improved to for any . Moreover, if the unique games conjecture is true then minimum vertex cover cannot be approximated within any constant factor better than 2. Although finding the minimum-size vertex cover is equivalent to finding the maximum-size independent set, as described above, the two problems are not equivalent in an approximation-preserving way: The Independent Set problem has no constant-factor approximation unless P = NP. Pseudocode APPROXIMATION-VERTEX-COVER(G) C = ∅ E'= G.E while E' ≠ ∅: let (u, v) be an arbitrary edge of E' C = C ∪ {u, v} remove from E' every edge incident on either u or v return C Applications Vertex cover optimization serves as a model for many real-world and theoretical problems. For example, a commercial establishment interested in installing the fewest possible closed circuit cameras covering all hallways (edges) connecting all rooms (nodes) on a floor might model the objective as a vertex cover minimization problem. The problem has also been used to model the elimination of repetitive DNA sequences for synthetic biology and metabolic engineering applications.	Mathematics	Graph theory	null
562883	https://en.wikipedia.org/wiki/Higher-order%20logic	Higher-order logic	In mathematics and logic, a higher-order logic (abbreviated HOL) is a form of logic that is distinguished from first-order logic by additional quantifiers and, sometimes, stronger semantics. Higher-order logics with their standard semantics are more expressive, but their model-theoretic properties are less well-behaved than those of first-order logic. The term "higher-order logic" is commonly used to mean higher-order simple predicate logic. Here "simple" indicates that the underlying type theory is the theory of simple types, also called the simple theory of types. Leon Chwistek and Frank P. Ramsey proposed this as a simplification of the complicated and clumsy ramified theory of types specified in the Principia Mathematica by Alfred North Whitehead and Bertrand Russell. Simple types is sometimes also meant to exclude polymorphic and dependent types. Quantification scope First-order logic quantifies only variables that range over individuals; second-order logic, also quantifies over sets; third-order logic also quantifies over sets of sets, and so on. Higher-order logic is the union of first-, second-, third-, ..., nth-order logic; i.e., higher-order logic admits quantification over sets that are nested arbitrarily deeply. Semantics There are two possible semantics for higher-order logic. In the standard or full semantics, quantifiers over higher-type objects range over all possible objects of that type. For example, a quantifier over sets of individuals ranges over the entire powerset of the set of individuals. Thus, in standard semantics, once the set of individuals is specified, this is enough to specify all the quantifiers. HOL with standard semantics is more expressive than first-order logic. For example, HOL admits categorical axiomatizations of the natural numbers, and of the real numbers, which are impossible with first-order logic. However, by a result of Kurt Gödel, HOL with standard semantics does not admit an effective, sound, and complete proof calculus. The model-theoretic properties of HOL with standard semantics are also more complex than those of first-order logic. For example, the Löwenheim number of second-order logic is already larger than the first measurable cardinal, if such a cardinal exists. The Löwenheim number of first-order logic, in contrast, is ℵ0, the smallest infinite cardinal. In Henkin semantics, a separate domain is included in each interpretation for each higher-order type. Thus, for example, quantifiers over sets of individuals may range over only a subset of the powerset of the set of individuals. HOL with these semantics is equivalent to many-sorted first-order logic, rather than being stronger than first-order logic. In particular, HOL with Henkin semantics has all the model-theoretic properties of first-order logic, and has a complete, sound, effective proof system inherited from first-order logic. Properties Higher-order logics include the offshoots of Church's simple theory of types and the various forms of intuitionistic type theory. Gérard Huet has shown that unifiability is undecidable in a type-theoretic flavor of third-order logic, that is, there can be no algorithm to decide whether an arbitrary equation between second-order (let alone arbitrary higher-order) terms has a solution. Up to a certain notion of isomorphism, the powerset operation is definable in second-order logic. Using this observation, Jaakko Hintikka established in 1955 that second-order logic can simulate higher-order logics in the sense that for every formula of a higher-order logic, one can find an equisatisfiable formula for it in second-order logic. The term "higher-order logic" is assumed in some context to refer to classical higher-order logic. However, modal higher-order logic has been studied as well. According to several logicians, Gödel's ontological proof is best studied (from a technical perspective) in such a context.	Mathematics	Mathematical logic	null
562958	https://en.wikipedia.org/wiki/Signs%20and%20symptoms	Signs and symptoms	Signs and symptoms are diagnostic indications of an illness, injury, or condition. Signs are objective and externally observable; symptoms are a person's reported subjective experiences. A sign for example may be a higher or lower temperature than normal, raised or lowered blood pressure or an abnormality showing on a medical scan. A symptom is something out of the ordinary that is experienced by an individual such as feeling feverish, a headache or other pains in the body, which occur as the body's immune system fights off an infection. Signs and symptoms Signs A medical sign is an objective observable indication of a disease, injury, or medical condition that may be detected during a physical examination. These signs may be visible, such as a rash or bruise, or otherwise detectable such as by using a stethoscope or taking blood pressure. Medical signs, along with symptoms, help in forming a diagnosis. Some examples of signs are nail clubbing of either the fingernails or toenails or an abnormal gait. Symptoms A symptom is something felt or experienced, such as pain or dizziness. Signs and symptoms are not mutually exclusive, for example a subjective feeling of fever can be noted as sign by using a thermometer that registers a high reading. The CDC lists various diseases by their signs and symptoms such as for measles which includes a high fever, conjunctivitis, and cough, followed a few days later by the measles rash. Cardinal signs and symptoms Cardinal signs and symptoms are very specific even to the point of being pathognomonic. A cardinal sign or cardinal symptom can also refer to the major sign or symptom of a disease. Abnormal reflexes can indicate problems with the nervous system. Signs and symptoms are also applied to physiological states outside the context of disease, as for example when referring to the signs and symptoms of pregnancy, or the symptoms of dehydration. Sometimes a disease may be present without showing any signs or symptoms when it is known as being asymptomatic. The disorder may be discovered through tests including scans. An infection may be asymptomatic but still be transmissible. Syndrome Signs and symptoms are often non-specific, but some combinations can be suggestive of certain diagnoses, helping to narrow down what may be wrong. A particular set of characteristic signs and symptoms that may be associated with a disorder is known as a syndrome. Related terms Symptomatic When a disease is evidenced by symptoms it is known as symptomatic. There are many conditions including subclinical infections that display no symptoms, and these are termed asymptomatic. Signs and symptoms may be mild or severe, brief or longer-lasting when they may become reduced (remission), or then recur (relapse or recrudescence) known as a flare-up. A flare-up may show more severe symptoms. The term chief complaint, also "presenting problem", is used to describe the initial concern of an individual when seeking medical help, and once this is clearly noted a history of the present illness may be taken. The symptom that ultimately leads to a diagnosis is called a cardinal symptom. Some symptoms can be misleading as a result of referred pain, where for example a pain in the right shoulder may be due to an inflamed gallbladder and not to presumed muscle strain. Prodrome Many diseases have an early prodromal stage where a few signs and symptoms may suggest the presence of a disorder before further specific symptoms may emerge. Measles for example has a prodromal presentation that includes a hacking cough, fever, and Koplik's spots in the mouth. Over half of migraine episodes have a prodromal phase. Schizophrenia has a notable prodromal stage, as has dementia. Nonspecific symptoms Some symptoms are specific, that is, they are associated with a single, specific medical condition. Nonspecific symptoms, sometimes also called equivocal symptoms, are not specific to a particular condition. They include unexplained weight loss, headache, pain, fatigue, loss of appetite, night sweats, and malaise. A group of three particular nonspecific symptoms – fever, night sweats, and weight loss – over a period of six months are termed B symptoms associated with lymphoma and indicate a poor prognosis. Other sub-types of symptoms include: constitutional or general symptoms, which affect general well-being or the whole body, such as a fever; concomitant symptoms, which are symptoms that occur at the same time as the primary symptom; prodromal symptoms, which are the first symptoms of an bigger set of problems; delayed symptoms, which happen some time after the trigger; and objective symptoms, which are symptoms whose existence can be observed and confirmed by a healthcare provider. Vital signs Vital signs are the four signs that can give an immediate measurement of the body's overall functioning and health status. They are temperature, heart rate, breathing rate, and blood pressure. The ranges of these measurements vary with age, weight, gender and with general health. A digital application has been developed for use in clinical settings that measures three of the vital signs (not temperature) using just a smartphone, and has been approved by NHS England. The application is registered as Lifelight First, and Lifelight Home is under development (2020) for monitoring-use by people at home using just the camera on their smartphone or tablet. This will additionally measure oxygen saturation and atrial fibrillation. Other devices are then not needed. Syndromes Many conditions are indicated by a group of known signs, or signs and symptoms. These can be a group of three known as a triad; a group of four ("tetrad"); or a group of five ("pentad"). An example of a triad is Meltzer's triad presenting purpura a rash, arthralgia painful joints, and myalgia painful and weak muscles. Meltzer's triad indicates the condition cryoglobulinemia. Huntington's disease is a neurodegenerative disease that is characterized by a triad of motor, cognitive, and psychiatric signs and symptoms. A large number of these groups that can be characteristic of a particular disease are known as a syndrome. Noonan syndrome for example, has a diagnostic set of unique facial and musculoskeletal features. Some syndromes such as nephrotic syndrome may have a number of underlying causes that are all related to diseases that affect the kidneys. Sometimes a child or young adult may have symptoms suggestive of a genetic disorder that cannot be identified even after genetic testing. In such cases the term SWAN (syndrome without a name) may be used. Often a diagnosis may be made at some future point when other more specific symptoms emerge but many cases may remain undiagnosed. The inability to diagnose may be due to a unique combination of symptoms or an overlap of conditions, or to the symptoms being atypical of a known disorder, or to the disorder being extremely rare. It is possible that a person with a particular syndrome might not display every single one of the signs and/or symptoms that compose/define a syndrome. Positive and negative Sensory symptoms can also be described as positive symptoms, or as negative symptoms depending on whether the symptom is abnormally present such as tingling or itchiness, or abnormally absent such as loss of smell. The following terms are used for negative symptoms – hypoesthesia is a partial loss of sensitivity to moderate stimuli, such as pressure, touch, warmth, cold. Anesthesia is the complete loss of sensitivity to stronger stimuli, such as pinprick. Hypoalgesia (analgesia) is loss of sensation to painful stimuli. Symptoms are also grouped in to negative and positive for some mental disorders such as schizophrenia. Positive symptoms are those that are present in the disorder and are not normally experienced by most individuals and reflects an excess or distortion of normal functions; examples are hallucinations, delusions, and bizarre behavior. Negative symptoms are functions that are normally found but that are diminished or absent, such as apathy and anhedonia. Dynamic and static Dynamic symptoms are capable of change depending on circumstance, whereas static symptoms are fixed or unchanging regardless of circumstance. For example, the symptoms of exercise intolerance are dynamic as they are brought on by exercise, but alleviate during rest. Fixed muscle weakness is a static symptom as the muscle will be weak regardless of exercise or rest. A majority of patients with metabolic myopathies have dynamic rather than static findings, typically experiencing exercise intolerance, muscle pain, and cramps with exercise rather than fixed weakness. Those with the metabolic myopathy of McArdle's disease (GSD-V) and some individuals with phosphoglucomutase deficiency (CDG1T/GSD-XIV), initially experience exercise intolerance during mild-moderate aerobic exercise, but the symptoms alleviate after 6–10 minutes in what is known as "second wind". Neuropsychiatric Neuropsychiatric symptoms are present in many degenerative disorders including dementia, and Parkinson's disease. Symptoms commonly include apathy, anxiety, and depression. Neurological and psychiatric symptoms are also present in some genetic disorders such as Wilson's disease. Symptoms of executive dysfunction are often found in many disorders including schizophrenia, and ADHD. Radiologic Radiologic signs are abnormal medical findings on imaging scanning. These include the Mickey Mouse sign and the Golden S sign. When using imaging to find the cause of a complaint, another unrelated finding may be found known as an incidental finding. Cardinal Cardinal signs and symptoms are those that may be diagnostic, and pathognomonic – of a certainty of diagnosis. Inflammation for example has a recognised group of cardinal signs and symptoms, as does exacerbations of chronic bronchitis, and Parkinson's disease. In contrast to a pathognomonic cardinal sign, the absence of a sign or symptom can often rule out a condition. This is known by the Latin term sine qua non. For example, the absence of known genetic mutations specific for a hereditary disease would rule out that disease. Another example is where the vaginal pH is less than 4.5, a diagnosis of bacterial vaginosis would be excluded. Reflexes A reflex is an automatic response in the body to a stimulus. Its absence, reduced (hypoactive), or exaggerated (hyperactive) response can be a sign of damage to the central nervous system or peripheral nervous system. In the patellar reflex (knee-jerk) for example, its reduction or absence is known as Westphal's sign and may indicate damage to lower motor neurons. When the response is exaggerated damage to the upper motor neurons may be indicated. Facies A number of medical conditions are associated with a distinctive facial expression or appearance known as a facies. An example is elfin facies which has facial features like those of the elf, and this may be associated with Williams syndrome, or Donohue syndrome. The most well-known facies is probably the Hippocratic facies that is seen on a person as they near death. Anamnestic signs Anamnestic signs (from anamnēstikós, ἀναμνηστικός, "able to recall to mind") are signs that indicate a past condition, for example paralysis in an arm may indicate a past stroke. Asymptomatic Some diseases including cancers, and infections may be present but show no signs or symptoms and these are known as asymptomatic. A gallstone may be asymptomatic and only discovered as an incidental finding. Easily spreadable viral infections such as COVID-19 may be asymptomatic but may still be transmissible. History Symptomatology A symptom (from Greek σύμπτωμα, "accident, misfortune, that which befalls", from συμπίπτω, "I befall", from συν- "together, with" and πίπτω, "I fall") is a departure from normal function or feeling. Symptomatology (also called semiology) is a branch of medicine dealing with the signs and symptoms of a disease. This study also includes the indications of a disease. It was first described as semiotics by Henry Stubbe in 1670 a term now used for the study of sign communication. Prior to the nineteenth century there was little difference in the powers of observation between physician and patient. Most medical practice was conducted as a co-operative interaction between the physician and patient; this was gradually replaced by a "monolithic consensus of opinion imposed from within the community of medical investigators". Whilst each noticed much the same things, the physician had a more informed interpretation of those things: "the physicians knew what the findings meant and the layman did not". Development of medical testing A number of advances introduced mostly in the 19th century, allowed for more objective assessment by the physician in search of a diagnosis, and less need of input from the patient. During the 20th century the introduction of a wide range of imaging techniques and other testing methods such as genetic testing, clinical chemistry tests, molecular diagnostics and pathogenomics have made a huge impact on diagnostic capability. In 1761 the percussion technique for diagnosing respiratory conditions was discovered by Leopold Auenbrugger. This method of tapping body cavities to note any abnormal sounds had already been in practice for a long time in cardiology. Percussion of the thorax became more widely known after 1808 with the translation of Auenbrugger's work from Latin into French by Jean-Nicolas Corvisart. In 1819 the introduction of the stethoscope by René Laennec began to replace the centuries-old technique of immediate auscultation – listening to the heart by placing the ear directly on the chest, with mediate auscultation using the stethoscope to listen to the sounds of the heart and respiratory tract. Laennec's publication was translated into English, 1824, by John Forbes. The 1846 introduction by surgeon John Hutchinson (1811–1861) of the spirometer, an apparatus for assessing the mechanical properties of the lungs via measurements of forced exhalation and forced inhalation. (The recorded lung volumes and air flow rates are used to distinguish between restrictive disease (in which the lung volumes are decreased: e.g., cystic fibrosis) and obstructive diseases (in which the lung volume is normal but the air flow rate is impeded; e.g., emphysema).) The 1851 invention by Hermann von Helmholtz (1821–1894) of the ophthalmoscope, which allowed physicians to examine the inside of the human eye. The () immediate widespread clinical use of Sir Thomas Clifford Allbutt's (1836–1925) six-inch (rather than twelve-inch) pocket clinical thermometer, which he had devised in 1867. The 1882 introduction of bacterial cultures by Robert Koch, initially for tuberculosis, being the first laboratory test to confirm bacterial infections. The 1895 clinical use of X-rays which began almost immediately after they had been discovered that year by Wilhelm Conrad Röntgen (1845–1923). The 1896 introduction of the sphygmomanometer, designed by Scipione Riva-Rocci (1863–1937), to measure blood pressure. Diagnosis The recognition of signs, and noting of symptoms may lead to a diagnosis. Otherwise a physical examination may be carried out, and a medical history taken. Further diagnostic medical tests such as blood tests, scans, and biopsies, may be needed. An X-ray for example would soon be diagnostic of a suspected bone fracture. A noted significance detected during an examination or from a medical test may be known as a medical finding. Examples Ascites (build-up of fluid in the abdomen) Nail clubbing (deformed nails) Cough Death rattle (last moments of life) Hemoptysis (blood-stained sputum) Jaundice Organomegaly an enlarged organ such as the liver (hepatomegaly) Palmar erythema (reddening of hands) Hypersalivation excessive (saliva) Unintentional weight loss	Biology and health sciences	Symptoms and signs	Health
563137	https://en.wikipedia.org/wiki/Sun%20bear	Sun bear	The sun bear (Helarctos malayanus) is a bear species in the family Ursidae found in the tropical forests of Southeast Asia. It is the only species in the genus Helarctos and the smallest bear species, standing nearly at the shoulder and weighing . It is stockily built, with large paws, strongly curved claws, small, rounded ears and a short snout. The fur is generally short and jet black, but can vary from grey to red. The sun bear gets its name from its characteristic orange to cream-coloured chest patch. Its unique morphology—inward-turned front feet, flattened chest, powerful forelimbs with large claws—suggests adaptations for climbing. The most arboreal (tree-living) of all bears, the sun bear is an excellent climber and sunbathes or sleeps in trees above the ground. It is mainly active during the day, though nocturnality might be more common in areas frequented by humans. Sun bears tend to remain solitary, but sometimes occur in twos (such as a mother and her cub). They do not seem to hibernate, possibly because food resources are available the whole year throughout the range. Being omnivores, sun bears' diet includes ants, bees, beetles, honey, termites, and plant material such as seeds and several kinds of fruits; vertebrates such as birds and deer are also eaten occasionally. They breed throughout the year; individuals become sexually mature at two to four years of age. Litters comprise one or two cubs that remain with their mother for around three years. The range of the sun bear is bounded by northeastern India to the north then south to southeast through Bangladesh, Cambodia, Myanmar, Laos, Thailand, and Vietnam in mainland Asia to Brunei, Indonesia, and Malaysia to the south. These bears are threatened by heavy deforestation and illegal hunting for food and the wildlife trade; they are also harmed in conflicts with humans when they enter farmlands, plantations, and orchards. The global population is estimated to have declined by 35% since the 1990s. The IUCN has listed this species as vulnerable. Etymology The generic name Helarctos comes from two Greek words: (, related to the sun) and (, bear). Another name is honey bear, in Malay and Indonesian, in reference to its habit of feeding on honey from honeycombs. Taxonomy and phylogeny The scientific name Ursus malayanus was proposed by Stamford Raffles in 1821 when he first described a sun bear from Sumatra. In 1825, Thomas Horsfield placed the species in a genus of its own, Helarctos, when describing a sun bear from Borneo. Subspecies and distribution H. annamiticus, described by Pierre Marie Heude in 1901 from Annam, is not considered a distinct species, but is subordinated as a junior synonym to H. m. malayanus. In 1906, Richard Lydekker proposed another subspecies by the name H. m. wardii for a sun bear skull, noting its similarities to a skull from Tibet with a thicker coat, but the Tibetan specimen was later found to be an Asian black bear (Ursus thibetanus). Genetic differences between the two subspecies are obscure and some authorities consider the species monotypic. Phylogeny The phylogenetic relationships among ursid species have remained ambiguous over the years. Noting the production of fertile hybrids between sun bears and sloth bears (Melursus ursinus), it was proposed that Helarctos be treated as a synonym of Melursus. However, studies differed on whether the two species were closely related. The brown bear/polar bear genetic lineage was estimated to have genetically diverged from the two black bears/sun bear lineage around (mya); the sun bear appears to have diverged from the two black bears between 6.26 and 5.09 mya. and 5.89–3.51 mya. Nuclear gene sequencing of bear species revealed that the sloth bear and the sun bear were the first Ursinae bears that radiated and are not included in the monophyletic Ursus group; moreover, all relationships between the bears were well resolved. Characteristics The sun bear is named so for its characteristic orange- to cream-coloured, crescent-like chest patch. It is the smallest of all bear species. It is stockily built, with large paws, strongly curved claws, small rounded ears and a short snout. The head-and-body length is between , and the shoulder height is nearly . Adults weigh . The snout is grey, silver, or orange. The fur is generally jet black, but can vary from grey to red. The hair is silky and fine, and is the shortest of all bear species, suiting its hot tropical habitat. The characteristic chest patch, typically U-shaped, but sometimes circular or spotlike, varies from orange or ochre-yellow to buff or cream, or even white. Some individuals may even lack the patch. Sun bears can expose the patch while standing on their hind feet as a threat display against enemies. Infants are greyish black with a pale brown or white snout and the chest patch is dirty white; the coat of older juveniles may be dark brown. The underfur is particularly thick and black in adults, while the guard hairs are lighter. Two whorls occur on the shoulders, from whence the hair radiates in all directions. A crest is seen on the sides of the neck and a whorl occurs in the centre of the breast patch. The edges of the paws are tan or brown, and the soles are fur-less, which possibly is an adaptation for climbing trees. The claws are sickle-shaped; the front claws are long and heavy. The tail is long. The sympatric Asian black bear has cream-coloured chest markings of a similar shape as those of sun bears and different claw markings. During feeding, the sun bear can extend its exceptionally long tongue to extract insects and honey. The teeth are very large, especially the canines, and the bite force quotient is high relative to its body size for reasons not well understood; a possible explanation could be its frequent opening of tropical hardwood trees with its powerful jaws and claws in pursuit of insects, larvae, or honey. The bite force is high for its size: a 50 Kg sun bear bites with a maximum force of 1907.3–2020.6 Newtons on the rear molar. The head is large, broad and heavy in proportion to the body, but the ears are proportionately smaller; the palate is wide in proportion to the skull. The overall unique morphology of this bear, such as its inward-turned front feet, flattened chest, and powerful fore limbs with large claws, indicates adaptations for extensive climbing. Ecology and behaviour Sun bears lead the most arboreal (tree-living) lifestyle among all bears. They are mainly active during the day, although nocturnality might be more common in areas frequented by humans. The sun bear is an excellent climber; it sunbathes or sleeps in trees above the ground. Bedding sites consist mainly of fallen hollow logs, but they also rest in standing trees with cavities, in cavities underneath fallen logs or tree roots, and in tree branches high above the ground. It is also an efficient swimmer. Sun bears are noted for their intelligence; a captive bear observed sugar being stored in a cupboard locked by a key, and later used its claw to open the lock. A study published in 2019 described skillful mimicry of facial expressions by sun bears, with precision comparable to that seen in gorillas and humans. Sun bears are shy and reclusive animals, and usually do not attack humans unless provoked to do so, or if they are injured or with their cubs; their timid nature led these bears to be often tamed and kept as pets in the past. Other sources, though, state that sun bears are known as very fierce animals when surprised in the forest. They are typically solitary but are sometimes seen in pairs (such as mothers and cubs). Sun bears stand on their hind feet for a broader view of their surroundings or smell far-off objects; they try to intimidate their enemies by displaying their chest patch if threatened. Vocalisations include grunts and snuffles while foraging for insects, and roars similar to those of a male orangutan during the breeding season; less commonly, they may give out short barks (like a rhinoceros) when they are surprised. Sun bears do not seem to hibernate, possibly because food resources are available the whole year throughout the range. They occupy home ranges of varying sizes in different areas, ranging from in Borneo and peninsular Malaysia; and in Ulu Segama Forest Reserve in Sabah. Tigers are their major predators; dholes and leopards have also been recorded preying on sun bears, but cases are relatively few. In one incident, a tiger-sun bear interaction resulted in a prolonged altercation and in the death of both animals. In another incident, a wild female sun bear was swallowed by a large reticulated python in East Kalimantan. Diet Sun bears are omnivores and feed on a broad variety of items, such as ants, bees, beetles, honey, termites, and plant material such as seeds and several kinds of fruits. Vertebrates such as birds, deer, eggs, and reptiles may be eaten occasionally. They forage mostly at night. Sun bears tear open hollow trees with their long, sharp claws and teeth in search of wild bees and honey. They also break termite mounds and quickly lick and suck the contents, holding pieces of the broken mound with their front paws. They consume figs in large amounts and eat them whole. In a study in the forests of Kalimantan, the fruits of Moraceae, Burseraceae, and Myrtaceae species made up more than 50% of the fruit diet; in times of fruit scarcity, sun bears switched to a more insectivorous diet. A study in Central Borneo revealed that sun bears play an important role in the seed dispersal of Canarium pilosum (a tree in the family Burseraceae). Sun bears eat the centre of coconut palms, and crush oil-rich seeds such as acorns. Oil palms are nutritious but not enough for subsistence. Reproduction Sun bears are polyoestrous; births occur throughout the year. Oestrus lasts five to seven days. Sun bears become sexually mature at two to four years of age. Reported lengths for pregnancies vary from 95 to 240 days; pregnancy tends to be longer in zoos in temperate climate possibly due to delay in implantation or fertilisation. Births occur inside hollow tree cavities. A litter typically comprises one or two cubs weighing around each. Cubs are born deaf with eyes closed. The eyes open at nearly 25 days, but they remain blind till 50 days after birth; the sense of hearing improves over the first 50 days. Cubs younger than two months are dependent on external stimulation for defecation. Cubs are kept on buttress roots at the base of trees until they learn how to walk and climb properly. Mothers protect their cubs aggressively. Offspring remain with their mother for nearly the first three years of their lives. Lifespan in captivity is generally over 20 years; one individual has lived for 34 years. Distribution and habitat The sun bear is native to the tropical forests of Southeast Asia; its range is bound by northeastern India to the north and extends south to Bangladesh, Myanmar, Thailand, Cambodia, Laos, and Vietnam to Brunei, Indonesia, and Malaysia to the south. Its presence in China was confirmed in 2017 when it was sighted in Yingjiang County of Yunnan Province. It is extinct in Singapore. These bears dwell primarily in two main types of forests throughout their range - deciduous and seasonally evergreen forests to the north of the Isthmus of Kra, and nonseasonal evergreen forests in Indonesia and Malaysia. They are typically found at low altitudes, such as below in western Thailand and peninsular Malaysia, but this varies widely throughout the range; in India, larger numbers have been recorded at elevations up to than in low-lying areas, probably due to habitat loss at ground level. They occur in montane areas in northeast India, but may not extend farther north into the unfavourable and colder Himalayan region; their distribution might be restricted to the northwest due to competition with sloth bears. The sun bear is sympatric with the Asian black bear throughout the remaining areas in the mainland range featuring a mix of seasonal forest types, with monthly rainfall below for a long spell of 3–7 months. In mountainous areas, Asian black bears are more common than sun bears, probably due to scarcity of invertebrates on which to feed. The major habitats in southern Thailand and peninsular Malaysia are moist evergreen forests, with more or less unvarying climate and heavy rainfall throughout the year, and low-lying or montane dipterocarp forests. Mangroves may be inhabited, but usually only when they are close to preferred habitat types. The sun bear tends to avoid heavily logged forests and areas close to human settlement. However, sun bears have been seen in farmlands, plantations and orchards, where they may be considered vermin. A survey in Lower Kinabatangan Segama Wetlands showed that sun bears were feared but were not common in oil palm plantations; Bornean bearded pigs, elephants and macaques were far more damaging to crops. Sun bears have been reported preying on poultry and livestock. Fossil remains suggest its occurrence farther north during the Pleistocene; it may have occurred as far south as Java in the middle to Late Pleistocene. Fossils also known from the Middle Pleistocene of Thailand along with Stegodon, gaur, wild water buffalo, and other living and extinct mammals. Today, it has been eliminated from the majority of its erstwhile range, especially in Thailand; populations are declining in most of the range countries. It disappeared from Singapore during the 1800s and 1900s, possibly due to extensive deforestation. Sun bear populations appear to decrease in size northward from Sundaland, and numbers are especially low in the northern and western extremes of the range. This has possibly been the case since prehistoric times, and is not a result of human interference. The population density varies from in Khao Yai National Park to in the Harapan Rainforest in southern Sumatra. Threats According to the IUCN Bear Specialist Group, sun bear populations have fallen by an estimated 35% since the 1990s. Numbers are especially low in Bangladesh and China, and populations in Vietnam are feared to decline severely by 50–80% in the next 30 years. Habitat fragmentation is on the rise, particularly in Borneo, Sumatra, and some areas of the mainland range. Heavy deforestation (due to agriculture, logging, and forest fires) and hunting for wildlife trade are severe threats throughout the range; human-bear conflicts are a relatively minor threat. Compared to other continents, Southeast Asia has undergone severe depletion in forest cover over the past few decades (by almost 12% between 1990 and 2010); this has resulted in substantial habitat loss for forest-dependent species such as sun bears. A 2007 study in East Borneo recorded severe loss of habitat and food resources due to droughts and forest fires brought about by the El Niño. With lack of research in predation, sources have documented very few predation events. In the island of Borneo sun bears were found to be hunted by python in their most vulnerable state. Pythons are successfully able to attack by taking advantage of the nighttime when the sun bears are sleep or nursing their cub. In Southeast Asia, a male leopard (Panthera pardus) has been photographed with a sun bear cub being held by the throat. This reported case has been reported to be the second confirmed predator as of 2019. During surveys in Kalimantan between 1994 and 1997, interviewees admitted to hunting sun bears and indicated that sun bear meat is eaten by indigenous people in several areas there. Studies have found evidence of pet trade and sale of sun bear parts such as gall bladders in traditional Chinese medicine (TCM) shops in Sabah and Sarawak. In 2018 and 2019, 128 TCM outlets in 24 locations across Sabah and Sarawak were surveyed and bear parts and derivatives were recorded for sale in 25% of the outlets surveyed, many of which would have been derived from locally sourced sun bears. Sun bears were killed by shooting or administering poison to protect coconut and snakefruit plantations in east Kalimantan. A report published by TRAFFIC in 2011 showed that sun bears, along with Asian black bears and brown bears, are specifically targeted for the bear bile trade in Southeast Asia, and are kept in bear farms in Laos, Vietnam, and Myanmar. Poaching is common in several countries in the region. Hunting pressure is rising even in some protected areas; in the Nam Ha National Protected Area in Laos, hunter snares have been found that specifically target bears. A study in Nagaland (northeastern India) recorded a sparse distribution of sun bears in the Fakim and Ntangki National Parks, and reported extensive illegal hunting for food and trade in bear parts. Protective laws have shown little success in controlling these threats, especially due to poor execution and high potential for gains by the trade. Conservation measures The sun bear is listed as vulnerable on the IUCN Red List, and is included in CITES Appendix I. With the exception of Sarawak (Malaysia) and Cambodia, the sun bear is legally protected from hunting in its whole range. A 2014 report documented rampant poaching and trade in sun bear parts in Sarawak, more than anywhere else in Malaysia; the researchers recommended stricter legislations in the state to protect local sun bears. The Bornean Sun Bear Conservation Centre, founded by Wong Siew Te in Sabah (Malaysia) in 2008, aims to work for the welfare of sun bears rescued from poor conditions in captivity and spread awareness about their conservation. The Malayan sun bears are part of an international captive-breeding program and a species survival plan under the Association of Zoos and Aquariums since late 1994. Since that same year, the European breed registry for sun bears is kept in the Cologne Zoological Garden, Germany.	Biology and health sciences	Bears	Animals
563161	https://en.wikipedia.org/wiki/Membrane%20potential	Membrane potential	Membrane potential (also transmembrane potential or membrane voltage) is the difference in electric potential between the interior and the exterior of a biological cell. It equals the interior potential minus the exterior potential. This is the energy (i.e. work) per charge which is required to move a (very small) positive charge at constant velocity across the cell membrane from the exterior to the interior. (If the charge is allowed to change velocity, the change of kinetic energy and production of radiation must be taken into account.) Typical values of membrane potential, normally given in units of milli volts and denoted as mV, range from –80 mV to –40 mV. For such typical negative membrane potentials, positive work is required to move a positive charge from the interior to the exterior. However, thermal kinetic energy allows ions to overcome the potential difference. For a selectively permeable membrane, this permits a net flow against the gradient. This is a kind of osmosis. Description All animal cells are surrounded by a membrane composed of a lipid bilayer with proteins embedded in it. The membrane serves as both an insulator and a diffusion barrier to the movement of ions. Transmembrane proteins, also known as ion transporter or ion pump proteins, actively push ions across the membrane and establish concentration gradients across the membrane, and ion channels allow ions to move across the membrane down those concentration gradients. Ion pumps and ion channels are electrically equivalent to a set of batteries and resistors inserted in the membrane, and therefore create a voltage between the two sides of the membrane. All plasma membranes have an electrical potential across them, with the inside usually negative with respect to the outside. The membrane potential has two basic functions. First, it allows a cell to function as a battery, providing power to operate a variety of "molecular devices" embedded in the membrane. Second, in electrically excitable cells such as neurons and muscle cells, it is used for transmitting signals between different parts of a cell. Signals in neurons and muscle cells Signals are generated in excitable cells by opening or closing of ion channels at one point in the membrane, producing a local change in the membrane potential. This change in the electric field can be quickly sensed by either adjacent or more distant ion channels in the membrane. Those ion channels can then open or close as a result of the potential change, reproducing the signal. In non-excitable cells, and in excitable cells in their baseline states, the membrane potential is held at a relatively stable value, called the resting potential. For neurons, resting potential is defined as ranging from –80 to –70 millivolts; that is, the interior of a cell has a negative baseline voltage of a bit less than one-tenth of a volt. The opening and closing of ion channels can induce a departure from the resting potential. This is called a depolarization if the interior voltage becomes less negative (say from –70 mV to –60 mV), or a hyperpolarization if the interior voltage becomes more negative (say from –70 mV to –80 mV). In excitable cells, a sufficiently large depolarization can evoke an action potential, in which the membrane potential changes rapidly and significantly for a short time (on the order of 1 to 100 milliseconds), often reversing its polarity. Action potentials are generated by the activation of certain voltage-gated ion channels. In neurons, the factors that influence the membrane potential are diverse. They include numerous types of ion channels, some of which are chemically gated and some of which are voltage-gated. Because voltage-gated ion channels are controlled by the membrane potential, while the membrane potential itself is influenced by these same ion channels, feedback loops that allow for complex temporal dynamics arise, including oscillations and regenerative events such as action potentials. Ion concentration gradients Differences in the concentrations of ions on opposite sides of a cellular membrane lead to a voltage called the membrane potential. Many ions have a concentration gradient across the membrane, including potassium (K+), which is at a high concentration inside and a low concentration outside the membrane. Sodium (Na+) and chloride (Cl−) ions are at high concentrations in the extracellular region, and low concentrations in the intracellular regions. These concentration gradients provide the potential energy to drive the formation of the membrane potential. This voltage is established when the membrane has permeability to one or more ions. In the simplest case, illustrated in the top diagram ("Ion concentration gradients"), if the membrane is selectively permeable to potassium, these positively charged ions can diffuse down the concentration gradient to the outside of the cell, leaving behind uncompensated negative charges. This separation of charges is what causes the membrane potential. The system as a whole is electro-neutral. The uncompensated positive charges outside the cell, and the uncompensated negative charges inside the cell, physically line up on the membrane surface and attract each other across the lipid bilayer. Thus, the membrane potential is physically located only in the immediate vicinity of the membrane. It is the separation of these charges across the membrane that is the basis of the membrane voltage. The top diagram is only an approximation of the ionic contributions to the membrane potential. Other ions including sodium, chloride, calcium, and others play a more minor role, even though they have strong concentration gradients, because they have more limited permeability than potassium. Physical basis The membrane potential in a cell derives ultimately from two factors: electrical force and diffusion. Electrical force arises from the mutual attraction between particles with opposite electrical charges (positive and negative) and the mutual repulsion between particles with the same type of charge (both positive or both negative). Diffusion arises from the statistical tendency of particles to redistribute from regions where they are highly concentrated to regions where the concentration is low. Voltage Voltage, which is synonymous with difference in electrical potential, is the ability to drive an electric current across a resistance. Indeed, the simplest definition of a voltage is given by Ohm's law: V=IR, where V is voltage, I is current and R is resistance. If a voltage source such as a battery is placed in an electrical circuit, the higher the voltage of the source the greater the amount of current that it will drive across the available resistance. The functional significance of voltage lies only in potential differences between two points in a circuit. The idea of a voltage at a single point is meaningless. It is conventional in electronics to assign a voltage of zero to some arbitrarily chosen element of the circuit, and then assign voltages for other elements measured relative to that zero point. There is no significance in which element is chosen as the zero point—the function of a circuit depends only on the differences not on voltages per se. However, in most cases and by convention, the zero level is most often assigned to the portion of a circuit that is in contact with ground. The same principle applies to voltage in cell biology. In electrically active tissue, the potential difference between any two points can be measured by inserting an electrode at each point, for example one inside and one outside the cell, and connecting both electrodes to the leads of what is in essence a specialized voltmeter. By convention, the zero potential value is assigned to the outside of the cell and the sign of the potential difference between the outside and the inside is determined by the potential of the inside relative to the outside zero. In mathematical terms, the definition of voltage begins with the concept of an electric field , a vector field assigning a magnitude and direction to each point in space. In many situations, the electric field is a conservative field, which means that it can be expressed as the gradient of a scalar function , that is, . This scalar field is referred to as the voltage distribution. The definition allows for an arbitrary constant of integration—this is why absolute values of voltage are not meaningful. In general, electric fields can be treated as conservative only if magnetic fields do not significantly influence them, but this condition usually applies well to biological tissue. Because the electric field is the gradient of the voltage distribution, rapid changes in voltage within a small region imply a strong electric field; on the converse, if the voltage remains approximately the same over a large region, the electric fields in that region must be weak. A strong electric field, equivalent to a strong voltage gradient, implies that a strong force is exerted on any charged particles that lie within the region. Ions and the forces driving their motion Electrical signals within biological organisms are, in general, driven by ions. The most important cations for the action potential are sodium (Na+) and potassium (K+). Both of these are monovalent cations that carry a single positive charge. Action potentials can also involve calcium (Ca2+), which is a divalent cation that carries a double positive charge. The chloride anion (Cl−) plays a major role in the action potentials of some algae, but plays a negligible role in the action potentials of most animals. Ions cross the cell membrane under two influences: diffusion and electric fields. A simple example wherein two solutions—A and B—are separated by a porous barrier illustrates that diffusion will ensure that they will eventually mix into equal solutions. This mixing occurs because of the difference in their concentrations. The region with high concentration will diffuse out toward the region with low concentration. To extend the example, let solution A have 30 sodium ions and 30 chloride ions. Also, let solution B have only 20 sodium ions and 20 chloride ions. Assuming the barrier allows both types of ions to travel through it, then a steady state will be reached whereby both solutions have 25 sodium ions and 25 chloride ions. If, however, the porous barrier is selective to which ions are let through, then diffusion alone will not determine the resulting solution. Returning to the previous example, let's now construct a barrier that is permeable only to sodium ions. Now, only sodium is allowed to diffuse cross the barrier from its higher concentration in solution A to the lower concentration in solution B. This will result in a greater accumulation of sodium ions than chloride ions in solution B and a lesser number of sodium ions than chloride ions in solution A. This means that there is a net positive charge in solution B from the higher concentration of positively charged sodium ions than negatively charged chloride ions. Likewise, there is a net negative charge in solution A from the greater concentration of negative chloride ions than positive sodium ions. Since opposite charges attract and like charges repel, the ions are now also influenced by electrical fields as well as forces of diffusion. Therefore, positive sodium ions will be less likely to travel to the now-more-positive B solution and remain in the now-more-negative A solution. The point at which the forces of the electric fields completely counteract the force due to diffusion is called the equilibrium potential. At this point, the net flow of the specific ion (in this case sodium) is zero. Plasma membranes Every cell is enclosed in a plasma membrane, which has the structure of a lipid bilayer with many types of large molecules embedded in it. Because it is made of lipid molecules, the plasma membrane intrinsically has a high electrical resistivity, in other words a low intrinsic permeability to ions. However, some of the molecules embedded in the membrane are capable either of actively transporting ions from one side of the membrane to the other or of providing channels through which they can move. In electrical terminology, the plasma membrane functions as a combined resistor and capacitor. Resistance arises from the fact that the membrane impedes the movement of charges across it. Capacitance arises from the fact that the lipid bilayer is so thin that an accumulation of charged particles on one side gives rise to an electrical force that pulls oppositely charged particles toward the other side. The capacitance of the membrane is relatively unaffected by the molecules that are embedded in it, so it has a more or less invariant value estimated at 2 μF/cm2 (the total capacitance of a patch of membrane is proportional to its area). The conductance of a pure lipid bilayer is so low, on the other hand, that in biological situations it is always dominated by the conductance of alternative pathways provided by embedded molecules. Thus, the capacitance of the membrane is more or less fixed, but the resistance is highly variable. The thickness of a plasma membrane is estimated to be about 7-8 nanometers. Because the membrane is so thin, it does not take a very large transmembrane voltage to create a strong electric field within it. Typical membrane potentials in animal cells are on the order of 100 millivolts (that is, one tenth of a volt), but calculations show that this generates an electric field close to the maximum that the membrane can sustain—it has been calculated that a voltage difference much larger than 200 millivolts could cause dielectric breakdown, that is, arcing across the membrane. Facilitated diffusion and transport The resistance of a pure lipid bilayer to the passage of ions across it is very high, but structures embedded in the membrane can greatly enhance ion movement, either actively or passively, via mechanisms called facilitated transport and facilitated diffusion. The two types of structure that play the largest roles are ion channels and ion pumps, both usually formed from assemblages of protein molecules. Ion channels provide passageways through which ions can move. In most cases, an ion channel is permeable only to specific types of ions (for example, sodium and potassium but not chloride or calcium), and sometimes the permeability varies depending on the direction of ion movement. Ion pumps, also known as ion transporters or carrier proteins, actively transport specific types of ions from one side of the membrane to the other, sometimes using energy derived from metabolic processes to do so. Ion pumps Ion pumps are integral membrane proteins that carry out active transport, i.e., use cellular energy (ATP) to "pump" the ions against their concentration gradient. Such ion pumps take in ions from one side of the membrane (decreasing its concentration there) and release them on the other side (increasing its concentration there). The ion pump most relevant to the action potential is the sodium–potassium pump, which transports three sodium ions out of the cell and two potassium ions in. As a consequence, the concentration of potassium ions K+ inside the neuron is roughly 30-fold larger than the outside concentration, whereas the sodium concentration outside is roughly five-fold larger than inside. In a similar manner, other ions have different concentrations inside and outside the neuron, such as calcium, chloride and magnesium. If the numbers of each type of ion were equal, the sodium–potassium pump would be electrically neutral, but, because of the three-for-two exchange, it gives a net movement of one positive charge from intracellular to extracellular for each cycle, thereby contributing to a positive voltage difference. The pump has three effects: (1) it makes the sodium concentration high in the extracellular space and low in the intracellular space; (2) it makes the potassium concentration high in the intracellular space and low in the extracellular space; (3) it gives the intracellular space a negative voltage with respect to the extracellular space. The sodium-potassium pump is relatively slow in operation. If a cell were initialized with equal concentrations of sodium and potassium everywhere, it would take hours for the pump to establish equilibrium. The pump operates constantly, but becomes progressively less efficient as the concentrations of sodium and potassium available for pumping are reduced. Ion pumps influence the action potential only by establishing the relative ratio of intracellular and extracellular ion concentrations. The action potential involves mainly the opening and closing of ion channels not ion pumps. If the ion pumps are turned off by removing their energy source, or by adding an inhibitor such as ouabain, the axon can still fire hundreds of thousands of action potentials before their amplitudes begin to decay significantly. In particular, ion pumps play no significant role in the repolarization of the membrane after an action potential. Another functionally important ion pump is the sodium-calcium exchanger. This pump operates in a conceptually similar way to the sodium-potassium pump, except that in each cycle it exchanges three Na+ from the extracellular space for one Ca++ from the intracellular space. Because the net flow of charge is inward, this pump runs "downhill", in effect, and therefore does not require any energy source except the membrane voltage. Its most important effect is to pump calcium outward—it also allows an inward flow of sodium, thereby counteracting the sodium-potassium pump, but, because overall sodium and potassium concentrations are much higher than calcium concentrations, this effect is relatively unimportant. The net result of the sodium-calcium exchanger is that in the resting state, intracellular calcium concentrations become very low. Ion channels Ion channels are integral membrane proteins with a pore through which ions can travel between extracellular space and cell interior. Most channels are specific (selective) for one ion; for example, most potassium channels are characterized by 1000:1 selectivity ratio for potassium over sodium, though potassium and sodium ions have the same charge and differ only slightly in their radius. The channel pore is typically so small that ions must pass through it in single-file order. Channel pores can be either open or closed for ion passage, although a number of channels demonstrate various sub-conductance levels. When a channel is open, ions permeate through the channel pore down the transmembrane concentration gradient for that particular ion. Rate of ionic flow through the channel, i.e. single-channel current amplitude, is determined by the maximum channel conductance and electrochemical driving force for that ion, which is the difference between the instantaneous value of the membrane potential and the value of the reversal potential. A channel may have several different states (corresponding to different conformations of the protein), but each such state is either open or closed. In general, closed states correspond either to a contraction of the pore—making it impassable to the ion—or to a separate part of the protein, stoppering the pore. For example, the voltage-dependent sodium channel undergoes inactivation, in which a portion of the protein swings into the pore, sealing it. This inactivation shuts off the sodium current and plays a critical role in the action potential. Ion channels can be classified by how they respond to their environment. For example, the ion channels involved in the action potential are voltage-sensitive channels; they open and close in response to the voltage across the membrane. Ligand-gated channels form another important class; these ion channels open and close in response to the binding of a ligand molecule, such as a neurotransmitter. Other ion channels open and close with mechanical forces. Still other ion channels—such as those of sensory neurons—open and close in response to other stimuli, such as light, temperature or pressure. Leakage channels Leakage channels are the simplest type of ion channel, in that their permeability is more or less constant. The types of leakage channels that have the greatest significance in neurons are potassium and chloride channels. Even these are not perfectly constant in their properties: First, most of them are voltage-dependent in the sense that they conduct better in one direction than the other (in other words, they are rectifiers); second, some of them are capable of being shut off by chemical ligands even though they do not require ligands in order to operate. Ligand-gated channels Ligand-gated ion channels are channels whose permeability is greatly increased when some type of chemical ligand binds to the protein structure. Animal cells contain hundreds, if not thousands, of types of these. A large subset function as neurotransmitter receptors—they occur at postsynaptic sites, and the chemical ligand that gates them is released by the presynaptic axon terminal. One example of this type is the AMPA receptor, a receptor for the neurotransmitter glutamate that when activated allows passage of sodium and potassium ions. Another example is the GABAA receptor, a receptor for the neurotransmitter GABA that when activated allows passage of chloride ions. Neurotransmitter receptors are activated by ligands that appear in the extracellular area, but there are other types of ligand-gated channels that are controlled by interactions on the intracellular side. Voltage-dependent channels Voltage-gated ion channels, also known as voltage dependent ion channels, are channels whose permeability is influenced by the membrane potential. They form another very large group, with each member having a particular ion selectivity and a particular voltage dependence. Many are also time-dependent—in other words, they do not respond immediately to a voltage change but only after a delay. One of the most important members of this group is a type of voltage-gated sodium channel that underlies action potentials—these are sometimes called Hodgkin-Huxley sodium channels because they were initially characterized by Alan Lloyd Hodgkin and Andrew Huxley in their Nobel Prize-winning studies of the physiology of the action potential. The channel is closed at the resting voltage level, but opens abruptly when the voltage exceeds a certain threshold, allowing a large influx of sodium ions that produces a very rapid change in the membrane potential. Recovery from an action potential is partly dependent on a type of voltage-gated potassium channel that is closed at the resting voltage level but opens as a consequence of the large voltage change produced during the action potential. Reversal potential The reversal potential (or equilibrium potential) of an ion is the value of transmembrane voltage at which diffusive and electrical forces counterbalance, so that there is no net ion flow across the membrane. This means that the transmembrane voltage exactly opposes the force of diffusion of the ion, such that the net current of the ion across the membrane is zero and unchanging. The reversal potential is important because it gives the voltage that acts on channels permeable to that ion—in other words, it gives the voltage that the ion concentration gradient generates when it acts as a battery. The equilibrium potential of a particular ion is usually designated by the notation Eion.The equilibrium potential for any ion can be calculated using the Nernst equation. For example, reversal potential for potassium ions will be as follows: where Eeq,K+= equilibrium potential for potassium, measured in volts R = universal gas constant, equal to 8.314 joules·K−1·mol−1 T = absolute temperature, measured in kelvins (= K = degrees Celsius + 273.15) z = number of elementary charges of the ion in question involved in the reaction F = Faraday constant, equal to 96,485 coulombs·mol−1 or J·V−1·mol−1 [K+]o= extracellular concentration of potassium, measured in mol·m−3 or mmol·l−1 [K+]i= intracellular concentration of potassium Even if two different ions have the same charge (i.e., K+ and Na+), they can still have very different equilibrium potentials, provided their outside and/or inside concentrations differ. Take, for example, the equilibrium potentials of potassium and sodium in neurons. The potassium equilibrium potential EK is −84 mV with 5 mM potassium outside and 140 mM inside. On the other hand, the sodium equilibrium potential, ENa, is approximately +66 mV with approximately 12 mM sodium inside and 140 mM outside. Changes to membrane potential during development A neuron's resting membrane potential actually changes during the development of an organism. In order for a neuron to eventually adopt its full adult function, its potential must be tightly regulated during development. As an organism progresses through development the resting membrane potential becomes more negative. Glial cells are also differentiating and proliferating as development progresses in the brain. The addition of these glial cells increases the organism's ability to regulate extracellular potassium. The drop in extracellular potassium can lead to a decrease in membrane potential of 35 mV. Cell excitability Cell excitability is the change in membrane potential that is necessary for cellular responses in various tissues. Cell excitability is a property that is induced during early embriogenesis. Excitability of a cell has also been defined as the ease with which a response may be triggered. The resting and threshold potentials forms the basis of cell excitability and these processes are fundamental for the generation of graded and action potentials. The most important regulators of cell excitability are the extracellular electrolyte concentrations (i.e. Na+, K+, Ca2+, Cl−, Mg2+) and associated proteins. Important proteins that regulate cell excitability are voltage-gated ion channels, ion transporters (e.g. Na+/K+-ATPase, magnesium transporters, acid–base transporters), membrane receptors and hyperpolarization-activated cyclic-nucleotide-gated channels. For example, potassium channels and calcium-sensing receptors are important regulators of excitability in neurons, cardiac myocytes and many other excitable cells like astrocytes. Calcium ion is also the most important second messenger in excitable cell signaling. Activation of synaptic receptors initiates long-lasting changes in neuronal excitability. Thyroid, adrenal and other hormones also regulate cell excitability, for example, progesterone and estrogen modulate myometrial smooth muscle cell excitability. Many cell types are considered to have an excitable membrane. Excitable cells are neurons, muscle (cardiac, skeletal, smooth), vascular endothelial cells, pericytes, juxtaglomerular cells, interstitial cells of Cajal, many types of epithelial cells (e.g. beta cells, alpha cells, delta cells, enteroendocrine cells, pulmonary neuroendocrine cells, pinealocytes), glial cells (e.g. astrocytes), mechanoreceptor cells (e.g. hair cells and Merkel cells), chemoreceptor cells (e.g. glomus cells, taste receptors), some plant cells and possibly immune cells. Astrocytes display a form of non-electrical excitability based on intracellular calcium variations related to the expression of several receptors through which they can detect the synaptic signal. In neurons, there are different membrane properties in some portions of the cell, for example, dendritic excitability endows neurons with the capacity for coincidence detection of spatially separated inputs. Equivalent circuit Electrophysiologists model the effects of ionic concentration differences, ion channels, and membrane capacitance in terms of an equivalent circuit, which is intended to represent the electrical properties of a small patch of membrane. The equivalent circuit consists of a capacitor in parallel with four pathways each consisting of a battery in series with a variable conductance. The capacitance is determined by the properties of the lipid bilayer, and is taken to be fixed. Each of the four parallel pathways comes from one of the principal ions, sodium, potassium, chloride, and calcium. The voltage of each ionic pathway is determined by the concentrations of the ion on each side of the membrane; see the Reversal potential section above. The conductance of each ionic pathway at any point in time is determined by the states of all the ion channels that are potentially permeable to that ion, including leakage channels, ligand-gated channels, and voltage-gated ion channels. For fixed ion concentrations and fixed values of ion channel conductance, the equivalent circuit can be further reduced, using the Goldman equation as described below, to a circuit containing a capacitance in parallel with a battery and conductance. In electrical terms, this is a type of RC circuit (resistance-capacitance circuit), and its electrical properties are very simple. Starting from any initial state, the current flowing across either the conductance or the capacitance decays with an exponential time course, with a time constant of , where is the capacitance of the membrane patch, and is the net resistance. For realistic situations, the time constant usually lies in the 1—100 millisecond range. In most cases, changes in the conductance of ion channels occur on a faster time scale, so an RC circuit is not a good approximation; however, the differential equation used to model a membrane patch is commonly a modified version of the RC circuit equation. Resting potential When the membrane potential of a cell goes for a long period of time without changing significantly, it is referred to as a resting potential or resting voltage. This term is used for the membrane potential of non-excitable cells, but also for the membrane potential of excitable cells in the absence of excitation. In excitable cells, the other possible states are graded membrane potentials (of variable amplitude), and action potentials, which are large, all-or-nothing rises in membrane potential that usually follow a fixed time course. Excitable cells include neurons, muscle cells, and some secretory cells in glands. Even in other types of cells, however, the membrane voltage can undergo changes in response to environmental or intracellular stimuli. For example, depolarization of the plasma membrane appears to be an important step in programmed cell death. The interactions that generate the resting potential are modeled by the Goldman equation. This is similar in form to the Nernst equation shown above, in that it is based on the charges of the ions in question, as well as the difference between their inside and outside concentrations. However, it also takes into consideration the relative permeability of the plasma membrane to each ion in question. The three ions that appear in this equation are potassium (K+), sodium (Na+), and chloride (Cl−). Calcium is omitted, but can be added to deal with situations in which it plays a significant role. Being an anion, the chloride terms are treated differently from the cation terms; the intracellular concentration is in the numerator, and the extracellular concentration in the denominator, which is reversed from the cation terms. Pi stands for the relative permeability of the ion type i. In essence, the Goldman formula expresses the membrane potential as a weighted average of the reversal potentials for the individual ion types, weighted by permeability. (Although the membrane potential changes about 100 mV during an action potential, the concentrations of ions inside and outside the cell do not change significantly. They remain close to their respective concentrations when then membrane is at resting potential.) In most animal cells, the permeability to potassium is much higher in the resting state than the permeability to sodium. As a consequence, the resting potential is usually close to the potassium reversal potential. The permeability to chloride can be high enough to be significant, but, unlike the other ions, chloride is not actively pumped, and therefore equilibrates at a reversal potential very close to the resting potential determined by the other ions. Values of resting membrane potential in most animal cells usually vary between the potassium reversal potential (usually around -80 mV) and around -40 mV. The resting potential in excitable cells (capable of producing action potentials) is usually near -60 mV—more depolarized voltages would lead to spontaneous generation of action potentials. Immature or undifferentiated cells show highly variable values of resting voltage, usually significantly more positive than in differentiated cells. In such cells, the resting potential value correlates with the degree of differentiation: undifferentiated cells in some cases may not show any transmembrane voltage difference at all. Maintenance of the resting potential can be metabolically costly for a cell because of its requirement for active pumping of ions to counteract losses due to leakage channels. The cost is highest when the cell function requires an especially depolarized value of membrane voltage. For example, the resting potential in daylight-adapted blowfly (Calliphora vicina) photoreceptors can be as high as -30 mV. This elevated membrane potential allows the cells to respond very rapidly to visual inputs; the cost is that maintenance of the resting potential may consume more than 20% of overall cellular ATP. On the other hand, the high resting potential in undifferentiated cells does not necessarily incur a high metabolic cost. This apparent paradox is resolved by examination of the origin of that resting potential. Little-differentiated cells are characterized by extremely high input resistance, which implies that few leakage channels are present at this stage of cell life. As an apparent result, potassium permeability becomes similar to that for sodium ions, which places resting potential in-between the reversal potentials for sodium and potassium as discussed above. The reduced leakage currents also mean there is little need for active pumping in order to compensate, therefore low metabolic cost. Graded potentials As explained above, the potential at any point in a cell's membrane is determined by the ion concentration differences between the intracellular and extracellular areas, and by the permeability of the membrane to each type of ion. The ion concentrations do not normally change very quickly (with the exception of Ca2+, where the baseline intracellular concentration is so low that even a small influx may increase it by orders of magnitude), but the permeabilities of the ions can change in a fraction of a millisecond, as a result of activation of ligand-gated ion channels. The change in membrane potential can be either large or small, depending on how many ion channels are activated and what type they are, and can be either long or short, depending on the lengths of time that the channels remain open. Changes of this type are referred to as graded potentials, in contrast to action potentials, which have a fixed amplitude and time course. As can be derived from the Goldman equation shown above, the effect of increasing the permeability of a membrane to a particular type of ion shifts the membrane potential toward the reversal potential for that ion. Thus, opening Na+ channels shifts the membrane potential toward the Na+ reversal potential, which is usually around +100 mV. Likewise, opening K+ channels shifts the membrane potential toward about –90 mV, and opening Cl− channels shifts it toward about –70 mV (resting potential of most membranes). Thus, Na+ channels shift the membrane potential in a positive direction, K+ channels shift it in a negative direction (except when the membrane is hyperpolarized to a value more negative than the K+ reversal potential), and Cl− channels tend to shift it towards the resting potential. Graded membrane potentials are particularly important in neurons, where they are produced by synapses—a temporary change in membrane potential produced by activation of a synapse by a single graded or action potential is called a postsynaptic potential. Neurotransmitters that act to open Na+ channels typically cause the membrane potential to become more positive, while neurotransmitters that activate K+ channels typically cause it to become more negative; those that inhibit these channels tend to have the opposite effect. Whether a postsynaptic potential is considered excitatory or inhibitory depends on the reversal potential for the ions of that current, and the threshold for the cell to fire an action potential (around –50mV). A postsynaptic current with a reversal potential above threshold, such as a typical Na+ current, is considered excitatory. A current with a reversal potential below threshold, such as a typical K+ current, is considered inhibitory. A current with a reversal potential above the resting potential, but below threshold, will not by itself elicit action potentials, but will produce subthreshold membrane potential oscillations. Thus, neurotransmitters that act to open Na+ channels produce excitatory postsynaptic potentials, or EPSPs, whereas neurotransmitters that act to open K+ or Cl− channels typically produce inhibitory postsynaptic potentials, or IPSPs. When multiple types of channels are open within the same time period, their postsynaptic potentials summate (are added together). Other values From the viewpoint of biophysics, the resting membrane potential is merely the membrane potential that results from the membrane permeabilities that predominate when the cell is resting. The above equation of weighted averages always applies, but the following approach may be more easily visualized. At any given moment, there are two factors for an ion that determine how much influence that ion will have over the membrane potential of a cell: That ion's driving force That ion's permeability If the driving force is high, then the ion is being "pushed" across the membrane. If the permeability is high, it will be easier for the ion to diffuse across the membrane. Driving force is the net electrical force available to move that ion across the membrane. It is calculated as the difference between the voltage that the ion "wants" to be at (its equilibrium potential) and the actual membrane potential (Em). So, in formal terms, the driving force for an ion = Em - Eion For example, at our earlier calculated resting potential of −73 mV, the driving force on potassium is 7 mV : (−73 mV) − (−80 mV) = 7 mV. The driving force on sodium would be (−73 mV) − (60 mV) = −133 mV. Permeability is a measure of how easily an ion can cross the membrane. It is normally measured as the (electrical) conductance and the unit, siemens, corresponds to 1 C·s−1·V−1, that is one coulomb per second per volt of potential. So, in a resting membrane, while the driving force for potassium is low, its permeability is very high. Sodium has a huge driving force but almost no resting permeability. In this case, potassium carries about 20 times more current than sodium, and thus has 20 times more influence over Em than does sodium. However, consider another case—the peak of the action potential. Here, permeability to Na is high and K permeability is relatively low. Thus, the membrane moves to near ENa and far from EK. The more ions are permeant the more complicated it becomes to predict the membrane potential. However, this can be done using the Goldman-Hodgkin-Katz equation or the weighted means equation. By plugging in the concentration gradients and the permeabilities of the ions at any instant in time, one can determine the membrane potential at that moment. What the GHK equations means is that, at any time, the value of the membrane potential will be a weighted average of the equilibrium potentials of all permeant ions. The "weighting" is the ions relative permeability across the membrane. Effects and implications While cells expend energy to transport ions and establish a transmembrane potential, they use this potential in turn to transport other ions and metabolites such as sugar. The transmembrane potential of the mitochondria drives the production of ATP, which is the common currency of biological energy. Cells may draw on the energy they store in the resting potential to drive action potentials or other forms of excitation. These changes in the membrane potential enable communication with other cells (as with action potentials) or initiate changes inside the cell, which happens in an egg when it is fertilized by a sperm. Changes in the dielectric properties of plasma membrane may act as hallmark of underlying conditions such as diabetes and dyslipidemia. In neuronal cells, an action potential begins with a rush of sodium ions into the cell through sodium channels, resulting in depolarization, while recovery involves an outward rush of potassium through potassium channels. Both of these fluxes occur by passive diffusion. A dose of salt may trigger the still-working neurons of a fresh cut of meat into firing, causing muscle spasms.	Biology and health sciences	Cell parts	Biology
563290	https://en.wikipedia.org/wiki/Adansonia	Adansonia	Adansonia is a genus made up of eight species of medium-to-large deciduous trees known as baobabs ( or ) or adansonias. They are placed in the family Malvaceae, subfamily Bombacoideae. They are native to Madagascar, mainland Africa, and Australia. The trees have also been introduced to other regions such as Asia. A genomic and ecological analysis has suggested that the genus is Madagascan in origin. The generic name honours Michel Adanson, the French naturalist and explorer who described Adansonia digitata. The baobab is also known as the "upside down tree", a name that originates from its appearance and several myths. They are among the most long-lived of vascular plants and have large flowers that are reproductive for a maximum of 15 hours. The flowers open around dusk, opening so quickly that movement can be detected by the naked eye, and are faded by the next morning. The fruits are large, oval to round and berry-like and hold kidney-shaped seeds in a dry, pulpy matrix. In the early 21st century, baobabs in southern Africa began to die off rapidly from a cause yet to be determined. It is unlikely that disease or pests would be able to kill many trees so rapidly, and some have speculated that the die-off is a result of dehydration. Description Baobabs are long-lived deciduous, small to large trees from tall with broad trunks and compact crowns. Young trees usually have slender, tapering trunks, often with a swollen base. Mature trees have massive trunks that are bottle-shaped or cylindrical and tapered from bottom to top. The trunk is made of fibrous wood arranged in concentric rings, although rings are not always formed annually and so cannot be used to determine the age of individual trees. Tree diameter fluctuates with rainfall so it is thought that water may be stored in the trunk. Baobab trees have two types of shoots—long, green vegetative ones, and stout, woody reproductive ones. Branches can be massive and spread out horizontally from the trunk or are ascending. Adansonia gregorii is generally the smallest of the baobabs, rarely getting to over tall and often with multiple trunks. Both A. rubrostipa and A. madagascariensis are small to large trees, from tall. The other baobabs grow from tall, with diameter trunks. A. digitata, however, often has massive single or multiple trunks of up to diameter. Leaves Leaves are palmately compound in mature trees, but seedlings and regenerating shoots may have simple leaves. The transition to compound leaves comes with age and may be gradual. Leaves have 5–11 leaflets, with the largest ones in the middle and may be stalkless or with short petioles. Leaflets may have toothed or smooth edges, and may be hairless or have simple-to-clumped hairs. Baobabs have stipules at the base of the leaves, but the stipules are soon shed in most species. Baobabs are deciduous, shedding leaves during the dry season. Flowers In most Adansonia species, the flowers are borne on short erect or spreading stalks in the axils of the leaves near the tips of reproductive shoots. Only A. digitata has flowers and fruits set on long, hanging stalks. There is usually only a single flower in an axil, but sometimes flowers occur in pairs. They are large, showy and strongly scented. They only open near dusk. Opening is rapid and movement of the flower parts is fast enough to be visible. Most Adansonia species are pollinated by bats. Flowers may remain attached to the trees for several days, but the reproductive phase is very short, with pollen shed during the first night and stigmas shriveled by the morning. The flower is made up of an outer 5-lobed calyx, and an inner ring of petals set around a fused tube of stamens. The outer lobes of the calyx are usually green (brown in A. grandidieri) and in bud are joined almost to the tip. As the flower opens, the calyx lobes split apart and become coiled or bent back (reflexed) at the base of the flower. The inner surface of the lobes are silky-hairy and cream, pink, or red. Sometimes the lobes do not separate cleanly, distorting the shape of the flower as they bend back. The calyx lobes remain fused at the base, leaving a feature (calyx tube) that has nectar-producing tissue and that is cup-shaped, flat or tubular; the form of the calyx tube varies with species. The flowers have a central tube (staminal tube) made up of fused stalks of stamens (filaments), with unfused filaments above. A densely hairy ovary is enclosed in the staminal tube, and a long style tipped with a stigma emerges from the filaments. Petals are set near the base of the staminal tube and are variable in shape and colour. The flowers, when fresh, may be white, cream, bright yellow or dark red, but fade quickly, often turning reddish when dried. Fruit The fruit of the baobabs is one of their distinguishing features. It is large, oval-to-round, and berry-like in most species (usually less than long in A. madagascariensis.). It has a dry, hard outer shell of variable thickness. In most species, the shell is indehiscent (does not break open easily). A. gibbosa is the only species with fruits that crack while still on the tree, which then tend to break open upon landing on the ground. Inside the outer shell, kidney-shaped seeds 10–15(−20) mm long are set in a dry pulp. Taxonomy The earliest written reports of baobab are from a 14th-century travelogue by the Arab traveler Ibn Battuta. The first botanical description was in the De medicina Aegyptiorum by Prospero Alpini (1592), looking at fruits that he observed in Egypt from an unknown source. They were called Bahobab, possibly from the Arabic abū ḥibāb meaning "many-seeded fruit". The French explorer and botanist Michel Adanson (1727–1806) observed a baobab tree in 1749 on the island of Sor in Senegal, and wrote the first detailed botanical description of the full tree, accompanied with illustrations. Recognizing the connection to the fruit described by Alpini he called the genus Baobab. Linnaeus later renamed the genus Adansonia, to honour Adanson, but use of baobab as one of the common names has persisted. The genus Adansonia is in the subfamily Bombacoideae, within the family Malvaceae in the order Malvales. The subfamily Bombacoideae was previously treated as the Bombacaceae family but it is no longer recognized at the rank of family by the Angiosperm Phylogeny Group I 1998, II 2003 or the Kubitzki system 2003. There are eight accepted species of Adansonia. A new species (Adansonia kilima Pettigrew, et al.), was described in 2012, found in high-elevation sites in eastern and southern Africa. This, however, is no longer recognized as a distinct species but considered a synonym of A. digitata. Some high-elevation trees in Tanzania show different genetics and morphology, but further study is needed to determine if recognition of them as a separate species is warranted. The genus Adansonia is further divided into three sections. Section Adansonia includes only A. digitata. This species has hanging flowers and fruit, set on long flowering stalks. This is the type species for the genus Adansonia. All species of Adansonia except A. digitata are diploid; A. digitata is tetraploid. Section Brevitubae includes A. grandidieri and A. suarexensis. These are species with flower buds that set on short pedicles and that are approximately twice as long as wide. The other species are all classified within the section Longitubae. They also have flowers/fruits set on short pedicels, but the flower buds are five or more times as long as wide. Species , there are eight recognized species of Adansonia, with six endemic to Madagascar, one native to mainland Africa and the Arabian Peninsula, and one native to Australia. The mainland African species (Adansonia digitata) also occurs on Madagascar, but it is not a native of that island. Baobabs were introduced in ancient times to south Asia and during the colonial era to the Caribbean. They are also present in the island nation of Cape Verde. A ninth species was described in 2012 (Adansonia kilima Pettigrew, et al.) but is no longer recognized as a distinct species. The African and Australian baobabs are similar in appearance, and the oldest splits within Adansonia are likely no older than 15 million years; thus, the Australian species represents a long-distance trans-oceanic dispersal event from Africa. The lineage leading to Adansonia was found to have diverged from its closest relatives in Bombacoideae like Ceiba /Chorisia at the end of the Eocene, during a time of abrupt global climate cooling and drying, while a divergence of this Adansonia+Ceiba /Chorisia clade from Pachira was found to be more ancient, dating to the middle Eocene. Habitat The Malagasy species are important components of the Madagascar dry deciduous forests. Within that biome, Adansonia madagascariensis and A. rubrostipa occur specifically in the Anjajavy Forest, sometimes growing out of the tsingy limestone itself. A. digitata has been called "a defining icon of African bushland". The tree also grows wild in Sudan in the regions of Darfur and the state of Kordofan. The locals call it "Gongolaze" and use its fruits as food and medicine and use the tree trunks as reservoirs to save water. Ecology Baobabs store water in the trunk (up to ) to endure harsh drought conditions. All occur in seasonally arid areas, and are deciduous, shedding their leaves during the dry season. Across Africa, the oldest and largest baobabs began to die in the early 21st century, likely from a combination of drought and rising temperatures. The trees appear to become parched, then become dehydrated and unable to support their massive trunks. Baobabs are important as nest sites for birds, in particular the mottled spinetail and four species of weaver. Notable trees Radiocarbon dating has provided data on a few individuals of A. digitata. The Panke baobab in Zimbabwe was some 2,450 years old when it died in 2011, making it the oldest angiosperm ever documented, and two other trees—Dorslandboom in Namibia and Glencoe in South Africa—were estimated to be approximately 2,000 years old. Another specimen known as Grootboom was dated and found to be at least 1,275 years old. The Glencoe Baobab, a specimen of A. digitata in Limpopo Province, South Africa, was considered to be the largest living individual, with a maximum circumference of and a diameter of about . The tree has since split into two parts, so the widest individual trunk may now be that of the Sunland Baobab, or Platland tree, also in South Africa. The diameter of this tree at ground level is and its circumference at breast height is . Two large baobabs growing in Tsimanampetsotse National Park were also studied using radiocarbon dating. One called Grandmother is made up of three fused trunks of different ages, with the oldest part of the tree an estimated 1,600 years old. The second, "polygamous baobab", has six fused stems, and is an estimated 1,000 years old. Culinary uses Leaves The tree's leaves may be eaten as a leaf vegetable. Fruit The white pith in the fruit of the Australian baobab (A. gregorii) tastes like sherbet. It has an acidic, tart, citrus flavor. It is a good source of vitamin C, potassium, carbohydrates, and phosphorus. The dried fruit powder of A. digitata, baobab powder, contains about 11% water, 80% carbohydrates (50% fiber), and modest levels of various nutrients, including riboflavin, calcium, magnesium, potassium, iron, and phytosterols, with low levels of protein and fats. Vitamin C content, described as variable in different samples, was in a range of per of dried powder. In 2008, baobab dried fruit pulp was authorized in the EU as a safe food ingredient, and later in the year was granted GRAS (generally recognized as safe) status in the United States. In Angola, the dry fruit of A. digitata is usually boiled, and the broth is used for juices or as the base for a type of ice cream known as gelado de múcua. In Zimbabwe, the fruit of A. digitata is eaten fresh or the crushed crumbly pulp is stirred into porridge and drinks. In Tanzania, the dry pulp of A. digitata is added to sugarcane to aid fermentation in brewing (beermaking). Seed The seeds of some species are a source of vegetable oil. The fruit pulp and seeds of A. grandidieri and A. za are eaten fresh. Other uses Some baobab species are sources of fiber, dye, and fuel. Indigenous Australians used the native species A. gregorii for several products, making string from the root fibers and decorative crafts from the fruits. Baobab oil from the seed is also used in cosmetics, particularly in moisturizers. In culture Baobab trees hold cultural and spiritual significance in many African societies. They are often the sites of communal gatherings, storytelling, and rituals. An unusual baobab was the namesake of Kukawa, formerly the capital of the Bornu Empire southwest of Lake Chad in Central Africa. In the novel The Little Prince, the titular character takes care to root out baobabs that try to grow on his tiny planet home. The fearsome, grasping baobab trees, researchers have contended, were meant to represent Nazism attempting to destroy the planet. Gallery	Biology and health sciences	Others	null
564276	https://en.wikipedia.org/wiki/Delimiter	Delimiter	A delimiter is a sequence of one or more characters for specifying the boundary between separate, independent regions in plain text, mathematical expressions or other data streams. An example of a delimiter is the comma character, which acts as a field delimiter in a sequence of comma-separated values. Another example of a delimiter is the time gap used to separate letters and words in the transmission of Morse code. In mathematics, delimiters are often used to specify the scope of an operation, and can occur both as isolated symbols (e.g., colon in "") and as a pair of opposing-looking symbols (e.g., angled brackets in ). Delimiters represent one of various means of specifying boundaries in a data stream. Declarative notation, for example, is an alternate method (without the use of delimiters) that uses a length field at the start of a data stream to specify the number of characters that the data stream contains. Overview Delimiters may be characterized as field and record delimiters, or as bracket delimiters. Field and record delimiters Field delimiters separate data fields. Record delimiters separate groups of fields. For example, the CSV format uses a comma as the delimiter between fields, and an end-of-line indicator as the delimiter between records: fname,lname,age,salary nancy,davolio,33,$30000 erin,borakova,28,$25250 tony,raphael,35,$28700 This specifies a simple flat-file database table using the CSV file format. Bracket delimiters Bracket delimiters, also called block delimiters, region delimiters, or balanced delimiters, mark both the start and end of a region of text. Common examples of bracket delimiters include: Conventions Historically, computing platforms have used certain delimiters by convention. The following tables depict a few examples for comparison. Programming languages (	Technology	Programming languages	null
564380	https://en.wikipedia.org/wiki/Expression%20vector	Expression vector	An expression vector, otherwise known as an expression construct, is usually a plasmid or virus designed for gene expression in cells. The vector is used to introduce a specific gene into a target cell, and can commandeer the cell's mechanism for protein synthesis to produce the protein encoded by the gene. Expression vectors are the basic tools in biotechnology for the production of proteins. The vector is engineered to contain regulatory sequences that act as enhancer and promoter regions and lead to efficient transcription of the gene carried on the expression vector. The goal of a well-designed expression vector is the efficient production of protein, and this may be achieved by the production of significant amount of stable messenger RNA, which can then be translated into protein. The expression of a protein may be tightly controlled, and the protein is only produced in significant quantity when necessary through the use of an inducer. In some systems, however, the protein may be expressed constitutively. Escherichia coli is commonly used as the host for protein production, but other cell types may also be used. An example of the use of expression vector is the production of insulin, which is used for medical treatments of diabetes. Elements An expression vector has features that any vector may have, such as an origin of replication, a selectable marker, and a suitable site for the insertion of a gene like the multiple cloning site. The cloned gene may be transferred from a specialized cloning vector to an expression vector, although it is possible to clone directly into an expression vector. The cloning process is normally performed in Escherichia coli. Vectors used for protein production in organisms other than E.coli may have, in addition to a suitable origin of replication for its propagation in E. coli, elements that allow them to be maintained in another organism, and these vectors are called shuttle vectors. Elements for expression An expression vector must have elements necessary for gene expression. These may include a promoter, the correct translation initiation sequence such as a ribosomal binding site and start codon, a termination codon, and a transcription termination sequence. There are differences in the machinery for protein synthesis between prokaryotes and eukaryotes, therefore the expression vectors must have the elements for expression that are appropriate for the chosen host. For example, prokaryotes expression vectors would have a Shine-Dalgarno sequence at its translation initiation site for the binding of ribosomes, while eukaryotes expression vectors would contain the Kozak consensus sequence. The promoter initiates the transcription and is therefore the point of control for the expression of the cloned gene. The promoters used in expression vector are normally inducible, meaning that protein synthesis is only initiated when required by the introduction of an inducer such as IPTG. Gene expression however may also be constitutive (i.e. protein is constantly expressed) in some expression vectors. Low level of constitutive protein synthesis may occur even in expression vectors with tightly controlled promoters. Protein tags After the expression of the gene product, it may be necessary to purify the expressed protein; however, separating the protein of interest from the great majority of proteins of the host cell can be a protracted process. To make this purification process easier, a purification tag may be added to the cloned gene. This tag could be histidine (His) tag, other marker peptides, or a fusion partners such as glutathione S-transferase or maltose-binding protein. Some of these fusion partners may also help to increase the solubility of some expressed proteins. Other fusion proteins such as green fluorescent protein may act as a reporter gene for the identification of successful cloned genes, or they may be used to study protein expression in cellular imaging. Other Elements The expression vector is transformed or transfected into the host cell for protein synthesis. Some expression vectors may have elements for transformation or the insertion of DNA into the host chromosome, for example the vir genes for plant transformation, and integrase sites for chromosomal integration . Some vectors may include targeting sequence that may target the expressed protein to a specific location such as the periplasmic space of bacteria. Expression/Production systems Different organisms may be used to express a gene's target protein, and the expression vector used will therefore have elements specific for use in the particular organism. The most commonly used organism for protein production is the bacterium Escherichia coli. However, not all proteins can be successfully expressed in E. coli, or be expressed with the correct form of post-translational modifications such as glycosylations, and other systems may therefore be used. Bacterial The expression host of choice for the expression of many proteins is Escherichia coli as the production of heterologous protein in E. coli is relatively simple and convenient, as well as being rapid and cheap. A large number of E. coli expression plasmids are also available for a wide variety of needs. Other bacteria used for protein production include Bacillus subtilis. Most heterologous proteins are expressed in the cytoplasm of E. coli. However, not all proteins formed may be soluble in the cytoplasm, and incorrectly folded proteins formed in cytoplasm can form insoluble aggregates called inclusion bodies. Such insoluble proteins will require refolding, which can be an involved process and may not necessarily produce high yield. Proteins which have disulphide bonds are often not able to fold correctly due to the reducing environment in the cytoplasm which prevents such bond formation, and a possible solution is to target the protein to the periplasmic space by the use of an N-terminal signal sequence. Another possibility is to manipulate the redox environment of the cytoplasm. Other more sophisticated systems are also being developed; such systems may allow for the expression of proteins previously thought impossible in E. coli, such as glycosylated proteins. The promoters used for these vector are usually based on the promoter of the lac operon or the T7 promoter, and they are normally regulated by the lac operator. These promoters may also be hybrids of different promoters, for example, the Tac-Promoter is a hybrid of trp and lac promoters. Note that most commonly used lac or lac-derived promoters are based on the lacUV5 mutant which is insensitive to catabolite repression. This mutant allows for expression of protein under the control of the lac promoter when the growth medium contains glucose since glucose would inhibit gene expression if wild-type lac promoter is used. Presence of glucose nevertheless may still be used to reduce background expression through residual inhibition in some systems. Examples of E. coli expression vectors are the pGEX series of vectors where glutathione S-transferase is used as a fusion partner and gene expression is under the control of the tac promoter, and the pET series of vectors which uses a T7 promoter. It is possible to simultaneously express two or more different proteins in E. coli using different plasmids. However, when 2 or more plasmids are used, each plasmid needs to use a different antibiotic selection as well as a different origin of replication, otherwise one of the plasmids may not be stably maintained. Many commonly used plasmids are based on the ColE1 replicon and are therefore incompatible with each other; in order for a ColE1-based plasmid to coexist with another in the same cell, the other would need to be of a different replicon, e.g. a p15A replicon-based plasmid such as the pACYC series of plasmids. Another approach would be to use a single two-cistron vector or design the coding sequences in tandem as a bi- or poly-cistronic construct. Yeast A yeast commonly used for protein production is Pichia pastoris. Examples of yeast expression vector in Pichia are the pPIC series of vectors, and these vectors use the AOX1 promoter which is inducible with methanol. The plasmids may contain elements for insertion of foreign DNA into the yeast genome and signal sequence for the secretion of expressed protein. Proteins with disulphide bonds and glycosylation can be efficiently produced in yeast. Another yeast used for protein production is Kluyveromyces lactis and the gene is expressed, driven by a variant of the strong lactase LAC4 promoter. Saccharomyces cerevisiae is particularly widely used for gene expression studies in yeast, for example in yeast two-hybrid system for the study of protein-protein interaction. The vectors used in yeast two-hybrid system contain fusion partners for two cloned genes that allow the transcription of a reporter gene when there is interaction between the two proteins expressed from the cloned genes. Baculovirus Baculovirus, a rod-shaped virus which infects insect cells, is used as the expression vector in this system. Insect cell lines derived from Lepidopterans (moths and butterflies), such as Spodoptera frugiperda, are used as host. A cell line derived from the cabbage looper is of particular interest, as it has been developed to grow fast and without the expensive serum normally needed to boost cell growth. The shuttle vector is called bacmid, and gene expression is under the control of a strong promoter pPolh. Baculovirus has also been used with mammalian cell lines in the BacMam system. Baculovirus is normally used for production of glycoproteins, although the glycosylations may be different from those found in vertebrates. In general, it is safer to use than mammalian virus as it has a limited host range and does not infect vertebrates without modifications. Plant Many plant expression vectors are based on the Ti plasmid of Agrobacterium tumefaciens. In these expression vectors, DNA to be inserted into plant is cloned into the T-DNA, a stretch of DNA flanked by a 25-bp direct repeat sequence at either end, and which can integrate into the plant genome. The T-DNA also contains the selectable marker. The Agrobacterium provides a mechanism for transformation, integration of into the plant genome, and the promoters for its vir genes may also be used for the cloned genes. Concerns over the transfer of bacterial or viral genetic material into the plant however have led to the development of vectors called intragenic vectors whereby functional equivalents of plant genome are used so that there is no transfer of genetic material from an alien species into the plant. Plant viruses may be used as vectors since the Agrobacterium method does not work for all plants. Examples of plant virus used are the tobacco mosaic virus (TMV), potato virus X, and cowpea mosaic virus. The protein may be expressed as a fusion to the coat protein of the virus and is displayed on the surface of assembled viral particles, or as an unfused protein that accumulates within the plant. Expression in plant using plant vectors is often constitutive, and a commonly used constitutive promoter in plant expression vectors is the cauliflower mosaic virus (CaMV) 35S promoter. Mammalian Mammalian expression vectors offer considerable advantages for the expression of mammalian proteins over bacterial expression systems - proper folding, post-translational modifications, and relevant enzymatic activity. It may also be more desirable than other eukaryotic non-mammalian systems whereby the proteins expressed may not contain the correct glycosylations. It is of particular use in producing membrane-associating proteins that require chaperones for proper folding and stability as well as containing numerous post-translational modifications. The downside, however, is the low yield of product in comparison to prokaryotic vectors as well as the costly nature of the techniques involved. Its complicated technology, and potential contamination with animal viruses of mammalian cell expression have also placed a constraint on its use in large-scale industrial production. Cultured mammalian cell lines such as the Chinese hamster ovary (CHO), COS, including human cell lines such as HEK and HeLa may be used to produce protein. Vectors are transfected into the cells and the DNA may be integrated into the genome by homologous recombination in the case of stable transfection, or the cells may be transiently transfected. Examples of mammalian expression vectors include the adenoviral vectors, the pSV and the pCMV series of plasmid vectors, vaccinia and retroviral vectors, as well as baculovirus. The promoters for cytomegalovirus (CMV) and SV40 are commonly used in mammalian expression vectors to drive gene expression. Non-viral promoter, such as the elongation factor (EF)-1 promoter, is also known. Cell-free systems E. coli cell lysate containing the cellular components required for transcription and translation are used in this in vitro method of protein production. The advantage of such system is that protein may be produced much faster than those produced in vivo since it does not require time to culture the cells, but it is also more expensive. Vectors used for E. coli expression can be used in this system although specifically designed vectors for this system are also available. Eukaryotic cell extracts may also be used in other cell-free systems, for example, the wheat germ cell-free expression systems. Mammalian cell-free systems have also been produced. Applications Laboratory use Expression vector in an expression host is now the usual method used in laboratories to produce proteins for research. Most proteins are produced in E. coli, but for glycosylated proteins and those with disulphide bonds, yeast, baculovirus and mammalian systems may be used. Production of peptide and protein pharmaceuticals Most protein pharmaceuticals are now produced through recombinant DNA technology using expression vectors. These peptide and protein pharmaceuticals may be hormones, vaccines, antibiotics, antibodies, and enzymes. The first human recombinant protein used for disease management, insulin, was introduced in 1982. Biotechnology allows these peptide and protein pharmaceuticals, some of which were previously rare or difficult to obtain, to be produced in large quantity. It also reduces the risks of contaminants such as host viruses, toxins and prions. Examples from the past include prion contamination in growth hormone extracted from pituitary glands harvested from human cadavers, which caused Creutzfeldt–Jakob disease in patients receiving treatment for dwarfism, and viral contaminants in clotting factor VIII isolated from human blood that resulted in the transmission of viral diseases such as hepatitis and AIDS. Such risk is reduced or removed completely when the proteins are produced in non-human host cells. Transgenic plant and animals In recent years, expression vectors have been used to introduce specific genes into plants and animals to produce transgenic organisms, for example in agriculture it is used to produce transgenic plants. Expression vectors have been used to introduce a vitamin A precursor, beta-carotene, into rice plants. This product is called golden rice. This process has also been used to introduce a gene into plants that produces an insecticide, called Bacillus thuringiensis toxin or Bt toxin which reduces the need for farmers to apply insecticides since it is produced by the modified organism. In addition expression vectors are used to extend the ripeness of tomatoes by altering the plant so that it produces less of the chemical that causes the tomatoes to rot. There have been controversies over using expression vectors to modify crops due to the fact that there might be unknown health risks, possibilities of companies patenting certain genetically modified food crops, and ethical concerns. Nevertheless, this technique is still being used and heavily researched. Transgenic animals have also been produced to study animal biochemical processes and human diseases, or used to produce pharmaceuticals and other proteins. They may also be engineered to have advantageous or useful traits. Green fluorescent protein is sometimes used as tags which results in animal that can fluoresce, and this have been exploited commercially to produce the fluorescent GloFish. Gene therapy Gene therapy is a promising treatment for a number of diseases where a "normal" gene carried by the vector is inserted into the genome, to replace an "abnormal" gene or supplement the expression of particular gene. Viral vectors are generally used but other nonviral methods of delivery are being developed. The treatment is still a risky option due to the viral vector used which can cause ill-effects, for example giving rise to insertional mutation that can result in cancer. However, there have been promising results.	Biology and health sciences	Molecular biology	Biology
564384	https://en.wikipedia.org/wiki/Clock%20face	Clock face	A clock face is the part of an analog clock (or watch) that displays time through the use of a flat dial with reference marks, and revolving pointers turning on concentric shafts at the center, called hands. In its most basic, globally recognized form, the periphery of the dial is numbered 1 through 12 indicating the hours in a 12-hour cycle, and a short hour hand makes two revolutions in a day. A long minute hand makes one revolution every hour. The face may also include a second hand, which makes one revolution per minute. The term is less commonly used for the time display on digital clocks and watches. A second type of clock face is the 24-hour analog dial, widely used in military and other organizations that use 24-hour time. This is similar to the 12-hour dial above, except it has hours numbered 1–24 (or 0–23) around the outside, and the hour hand makes only one revolution per day. Some special-purpose clocks, such as timers and sporting event clocks, are designed for measuring periods less than one hour. Clocks can indicate the hour with Roman numerals or Hindu–Arabic numerals, or with non-numeric indicator marks. The two numbering systems have also been used in combination, with the prior indicating the hour and the latter the minute. Longcase clocks (grandfather clocks) typically use Roman numerals for the hours. Clocks using only Arabic numerals first began to appear in the mid-18th century. The clock face is so familiar that the numbers are often omitted and replaced with unlabeled graduations (marks), particularly in the case of watches. Occasionally, markings of any sort are dispensed with, and the time is read by the angles of the hands. Reading a modern clock face Most modern clocks have the numbers 1 through 12 printed at equally spaced intervals around the periphery of the face with the 12 at the top, indicating the hour, and on many models, sixty dots or lines evenly spaced in a ring around the outside of the dial, indicating minutes and seconds. The time is read by observing the placement of several "hands", which emanate from the centre of the dial: A short, thick "hour" hand; A long, thinner "minute" hand; On some models, a very thin "second" or "sweep" hand All three hands continuously rotate around the dial in a clockwise direction – in the direction of increasing numbers. The second, or sweep, hand moves relatively quickly, taking a full minute (sixty seconds) to make a complete rotation from 12 to 12. For every rotation of the second hand, the minute hand will move from one minute mark to the next. The minute hand rotates more slowly around the dial. It takes one hour (sixty minutes) to make a complete rotation from 12 to 12. For every rotation of the minute hand, the hour hand will move from one hour mark to the next. The hour hand moves slowest of all, taking half a day (twelve hours) to make a complete rotation. It starts from "12" at midnight, makes one rotation until it is pointing at "12" again at noon, and then makes another rotation until it is pointing at "12" again at midnight of the next morning. Historical development The word clock derives from the medieval Latin word for "bell"; , and has cognates in many European languages. Clocks spread to England from the Low Countries, so the English word came from the Middle Low German and Middle Dutch Klocke. The first mechanical clocks, built in 13th-century Europe, were striking clocks: their purpose was to ring bells upon the canonical hours, to call the local community to prayer. These were tower clocks installed in bell towers in public places, to ensure that the bells were audible over a wide area. Soon after these first mechanical clocks were in place clockmakers realized that their wheels could be used to drive an indicator on a dial on the outside of the tower, where it could be widely seen, so the local population could tell the time between the hourly strikes. Before the late 14th century, a fixed hand (often a carving literally shaped like a hand) indicated the hour by pointing to numbers on a rotating dial; after this time, the current convention of a rotating hand on a fixed dial was adopted. Minute hands (so named because they indicated the small, or minute, divisions of the hour) only came into regular use around 1690, after the invention of the pendulum and anchor escapement increased the precision of time-telling enough to justify it. In some precision clocks, a third hand, which rotated once a minute, was added in a separate subdial. This was called the "second-minute" hand (because it measured the secondary minute divisions of the hour), which was shortened to "second" hand. The convention of the hands moving clockwise evolved in imitation of the sundial. In the Northern hemisphere, where the clock face originated, the shadow of the gnomon on a horizontal sundial moves clockwise during the day. French decimal time During the French Revolution in 1793, in connection with its Republican calendar, France attempted to introduce a decimal time system. This had 10 decimal hours in the day, 100 decimal minutes per hour, and 100 decimal seconds per minute. Therefore, the decimal hour was more than twice as long (144 min) as the present hour, the decimal minute was slightly longer than the present minute (86.4 seconds) and the decimal second was slightly shorter (0.864 sec) than the present second. Clocks were manufactured with this alternate face, usually combined with traditional hour markings. However, it did not catch on, and France discontinued the mandatory use of decimal time on 7 April 1795, although some French cities used decimal time until 1801. Stylistic development Until the last quarter of the 17th century, hour markings were etched into metal faces and the recesses filled with black wax. Subsequently, higher contrast and improved readability was achieved with white enamel plaques painted with black numbers. Initially, the numbers were printed on small, individual plaques mounted on a brass substructure. This was not a stylistic decision, rather enamel production technology had not yet achieved the ability to create large pieces of enamel. The "13-piece face" was an early attempt to create an entirely white enamel face. As the name suggests, it was composed of 13 enamel plaques: 12 numbered wedges fitted around a circle. The first single-piece enamel faces, not unlike those in production today, began to appear . It is customary for modern advertisements to display clocks and watches set to approximately 10:10 or 1:50, as this V-shaped arrangement roughly makes a smile, imitates a human figure with raised arms, and leaves the watch company's logo unobscured by the hands. In the 1970s, German designer Tian Harlan invented the Chromachron, a wristwatch with a clock face that has no dials but a disc with pie-shaped pattern rotating by the minute over color patterns representing both hours and minutes. Technological obsolescence In the 2010s, some United Kingdom schools started replacing analogue clocks in examination halls with digital clocks because an increasing number of pupils were unable to read analogue clocks. Smartphone and computer clocks are often digital rather than analogue, and proponents of replacing analogue clock faces argue that they have become technologically obsolete. However, reading analogue clocks is still part of American elementary school curricula; proponents of analogue clocks argue that their inclusion in the curriculum reinforces basic mathematical concepts that are taught in elementary school.	Technology	Clocks	null
564971	https://en.wikipedia.org/wiki/African%20buffalo	African buffalo	The African buffalo (Syncerus caffer) is a large sub-Saharan African bovine. There are five subspecies that are recognized as valid by most authorities: Syncerus caffer caffer, the Cape buffalo, is the nominotypical subspecies, as well as the largest, found in Southern and East Africa. S. c. nanus, the forest buffalo, is the smallest subspecies, common in forest areas of Central and West Africa S. c. brachyceros, the Sudan buffalo, a smaller version of the Cape buffalo, found in the drier, northern areas of Central and West Africa. S. c. aequinoctialis, the Nile Buffalo, sometimes considered identical to the Sudan buffalo, found in the drier, northern areas of East and Central Africa. S. c. mathewsi, the mountain buffalo, a disputed subspecies from the Virunga Mountains in Central Africa. The adult African buffalo's horns are its characteristic feature: they have fused bases, forming a continuous bone shield across the top of the head, referred to as a "boss". The African buffalo is more closely related to other buffalo species than it is to other bovids such as American bison or domestic cattle, with its closest living relative being the Asian water buffalo. Its unpredictable temperament may be part of the reason that the African buffalo has never been domesticated, which would also explain why the African buffalo has no domesticated descendants, unlike the wild yak and wild water buffalo which are the ancestors of the domestic yak and water buffalo. Natural predators of adult African buffaloes include lions, African wild dogs, spotted hyenas, and Nile crocodiles. As one of the Big Five game animals, the Cape buffalo is a sought-after trophy in hunting. Description The African buffalo is a very robust species. Its shoulder height can range from and its head-and-body length can range from . The tail can range from long. Compared with other large bovids, it has a long but stocky body (the body length can exceed the wild water buffalo, which is heavier and taller) and short but thickset legs, resulting in a relatively short standing height. Cape buffaloes weigh (males weigh about more than females). In comparison, African forest buffaloes, at , are only half that size. Its head is carried low; its top is located below the backline. The front hooves of the buffalo are wider than the rear, which is associated with the need to support the weight of the front part of the body, which is heavier and more powerful than the back. Savannah-type buffaloes have black or dark brown coats with age. Old bulls often have whitish circles around their eyes and on their face. Females tend to have more reddish coats. Forest-type buffaloes are 30–40% smaller, reddish brown in colour, with much more hair growth around the ears and with horns that curve back and slightly up. Calves of both types have red coats. A characteristic feature of the horns of adult male African buffalo (southern and eastern populations) is that the bases come very close together, forming a shield referred to as a "boss". From the base, the horns diverge downwards, then smoothly curve upwards and outwards and in some cases inwards and or backwards. In large bulls, the distance between the ends of the horns can reach upwards of one metre (the record being 64.5 inches 164 cm). The horns form fully when the animal reaches the age of 5 or 6 years old, but the bosses do not become "hard" until it reaches the age of 8 to 9 years old. In cows, the horns are, on average, 10–20% smaller, and they do not have a boss. Forest-type buffalo horns are smaller than those of the savanna-type buffaloes from Southern and East Africa, usually measuring less than , and are almost never fused. Unlike other large bovines, African buffalo have 52 chromosomes (for comparison, American bison and domestic cattle have 60). This means domestic cattle and bison are unable to create hybrid offspring with cape buffalo. Subspecies Ecology The African buffalo is one of the most successful grazers in Africa. It lives in savannas, swamps and floodplains, as well as mopane grasslands, and the forests of the major mountains of Africa. This buffalo prefers a habitat with dense cover, such as reeds and thickets, but can also be found in open woodland. While not particularly demanding in regard to habitat, they require water daily, and so they depend on perennial sources of water. Like the plains zebra, the buffalo can live on tall, coarse grasses. Herds of buffalo mow down grasses and make way for more selective grazers. When feeding, the buffalo makes use of its tongue and wide incisor row to eat grass more quickly than most other African herbivores. Buffaloes do not stay on trampled or depleted areas for long. Other than humans, African buffaloes have few predators and are capable of defending themselves against (and killing) lions. Lions kill and eat buffaloes regularly, and in some regions, the buffaloes are the lions' primary prey. It often takes several lions to bring down a single adult buffalo, and the entire pride may join in the hunt. However, several incidents have been reported in which lone adult male lions have successfully brought down adult buffaloes. On very rare occasions, buffaloes and white rhinos will fight over territory; due to the rhino's strength and size advantage, the rhino typically wins and the buffalo can die from injuries sustained during the encounter. Rhinos live solitary lives, whereas buffalo (excluding solitary adult bulls) primarily live social lives and thus they do not usually recognize each other as threats. Hippopotamuses and buffalo also do not normally interact, but if the buffalo provokes the hippo or makes it feel threatened, a fight can break out, but this is also rare. Adolescent bull African elephants may harass or kill Cape buffalo, either out of territorial aggression or while in musth; when they do this, the calves are most likely to be killed by the elephant attack as they are defenseless when facing an elephant alone, whereas adults will try to fight back and may survive (or succumb to injuries afterward). The average-sized Nile crocodile typically attacks only old solitary animals and young calves, though they can kill healthy adults. Exceptionally large, old male crocodiles may become semi-habitual predators of buffaloes. The cheetah, leopard, African wild dog and spotted hyena are normally a threat only to newborn calves, though larger clans of hyenas have been recorded killing cows (mainly pregnant ones) and, on rare occasions, full-grown bulls. Large packs of wild dogs have been observed to hunt calves and sick adults. Diseases The African buffalo is susceptible to many diseases, including those shared with domestic cattle, such as bovine tuberculosis, corridor disease, and foot-and-mouth disease. As with many diseases, these problems remain dormant within a population as long as the health of the animals is good. These diseases do, however, restrict the legal movements of the animals and fencing infected areas from unaffected areas is enforced. Some wardens and game managers have managed to protect and breed "disease-free" herds which become very valuable because they can be transported. Most well-known are Lindsay Hunt's efforts to source uninfected animals from the Kruger National Park in South Africa. Some disease-free buffaloes in South Africa have been sold to breeders for close to US$130,000. Social behavior Herd size is highly variable. The core of the herds is made up of related females, and their offspring, in an almost linear dominance hierarchy. The basic herds are surrounded by subherds of subordinate males, high-ranking males and females, and old or invalid animals. African buffaloes engage in several types of group behavior. Females appear to exhibit a sort of "voting behavior". During resting time, the females stand up, shuffle around, and sit back down again. They sit in the direction they think they should move. After an hour of more shuffling, the females travel in the direction they decide. This decision is communal and not based on hierarchy or dominance. When chased by predators, a herd sticks close together and makes it hard for the predators to pick off one member. Calves are gathered in the middle. A buffalo herd responds to the distress call of a threatened member and tries to rescue it. A calf's distress call gets the attention of not only the mother, but also the herd. Buffaloes engage in mobbing behavior when fighting off predators. They have been recorded killing lions and chasing lions up trees and keeping them there for two hours, after the lions have killed a member of their group. Lion cubs can get trampled and killed. In one videotaped instance, known as the Battle at Kruger, a calf survived an attack by both lions and a crocodile after intervention of the herd. Males have a linear dominance hierarchy based on age and size. Since a buffalo is safer when a herd is larger, dominant bulls may rely on subordinate bulls and sometimes tolerate their copulation. The young males keep their distance from the dominant bull, which is recognizable by the thickness of his horns. Adult bulls spar in play, dominance interactions, or actual fights. A bull approaches another, lowing, with his horns down, and waits for the other bull to do the same thing. When sparring, the bulls twist their horns from side to side. If the sparring is for play, the bull may rub his opponent's face and body during the sparring session. Actual fights are violent but rare and brief. Calves may also spar in play, but adult females rarely spar at all. During the dry season, males split from the herd and form bachelor groups. Two types of bachelor herds occur: ones made of males aged four to seven years and those of males 12 years or older. During the wet season, the younger bulls rejoin a herd to mate with the females. They stay with them throughout the season to protect the calves. Some older bulls cease to rejoin the herd, as they can no longer compete with the younger, more aggressive males. The old bachelors are called dagga boys ("mud covered"), and are considered the most dangerous to humans. Vocalizations African buffaloes make various vocalizations. Many calls are lower-pitched versions of those emitted by domestic cattle. They emit low-pitched, two- to four-second calls intermittently at three- to six-second intervals to signal the herd to move. To signal to the herd to change direction, leaders emit "gritty", "creaking gate" sounds. When moving to drinking places, some individuals make long "maaa" calls up to 20 times a minute. When being aggressive, they make explosive grunts that may last long or turn into a rumbling growl. Cows produce croaking calls when looking for their calves. Calves make a similar call of a higher pitch when in distress. When threatened by predators, they make drawn-out "waaaa" calls. Dominant individuals make calls to announce their presence and location. A version of the same call, but more intense, is emitted as a warning to an encroaching inferior. When grazing, they make various sounds, such as brief bellows, grunts, honks, and croaks. Reproduction Females reach sexual maturity at around five years of age while males are sexually matured at four to six. African buffaloes mate and give birth only during the rainy seasons. Birth peak takes place early in the season, while mating peaks later. A bull closely guards a cow that comes into heat, while keeping other bulls at bay. This is difficult, as cows are quite evasive and attract many males to the scene. By the time a cow is in full estrus, only the most dominant bull in the herd/subherd is there. Cows first calve at five years of age, after a gestation period of 11.5 months. Newborn calves remain hidden in vegetation for the first few weeks while being nursed occasionally by the mother before joining the main herd. Older calves are held in the centre of the herd for safety. The maternal bond between mother and calf lasts longer than in most bovids. That bonding ends when a new calf is born, and the mother then keeps her previous offspring at bay with horn jabs. Nevertheless, the yearling follows its mother for another year or so. Males leave their mothers when they are two years old and join the bachelor groups. Young calves, unusually for bovids, suckle from behind their mothers, pushing their heads between the mothers' legs. In the wild African buffaloes have an average lifespan of 11 years but they've been recorded to reach 22 years of age. In captivity they can live for a maximum of 29.5 years though they only live 16 years on average. Relationship with humans Status The current status of the African buffalo is dependent on the animal's value to both trophy hunters and tourists, paving the way for conservation efforts through anti-poaching patrols, village crop damage payouts, and CAMPFIRE payback programs to local areas. The African buffalo is listed as Near threatened by the IUCN, with a decreasing population of 400,000 individuals. While some populations (subspecies) are decreasing, others will remain unchanged in the long term if large, healthy populations continue to persist in a substantial number of national parks, equivalent reserves and hunting zones in southern and eastern Africa." In the most recent and available census data at continental scale, the total estimated numbers of the three savanna-type African buffalo subspecies (S. c. caffer, S. c. brachyceros and S. c. aequinoctialis) are at 513,000 individuals. In the past, numbers of African buffaloes suffered their most severe collapse during the great rinderpest epidemic of the 1890s, which, coupled with pleuro-pneumonia, caused mortalities as high as 95% among livestock and wild ungulates. Being a member of the big five game group, a term used to describe the five most dangerous animals to hunt, the Cape buffalo is a sought-after trophy, with some hunters paying over $10,000 for the opportunity to hunt one. The larger bulls are targeted for their trophy value, although in some areas, buffaloes are still hunted for meat. Attacks One of the "big five" African game, it is known as "the Black Death" or "the widowmaker", and is widely regarded as a very dangerous animal. African buffaloes are sometimes reported to kill more people in Africa than any other animal, although the same claim is also made of hippopotamuses and crocodiles. These numbers may be somewhat overestimated; for example, in the country of Mozambique, attacks, especially fatal ones, were much less frequent on humans than those by hippos, and especially, Nile crocodiles. In Uganda, on the other hand, large herbivores were found to attack more people on average than lions or leopards and have a higher rate of inflicting fatalities during attacks than the predators (the African buffalo, in particular, killing humans in 49.5% of attacks on them), but hippos and even elephants may still kill more people per annum than buffaloes. African buffaloes are notorious among big-game hunters as very dangerous animals, with wounded animals reported to ambush and attack pursuers. Domestication The Cape buffalo hybridized with Indian water buffalo to create the Jafarabadi buffalo breed.	Biology and health sciences	Artiodactyla	null
565530	https://en.wikipedia.org/wiki/Mesopelagic%20zone	Mesopelagic zone	The mesopelagic zone (Greek μέσον, middle), also known as the middle pelagic or twilight zone, is the part of the pelagic zone that lies between the photic epipelagic and the aphotic bathypelagic zones. It is defined by light, and begins at the depth where only 1% of incident light reaches and ends where there is no light; the depths of this zone are between approximately 200 to 1,000 meters (~656 to 3,280 feet) below the ocean surface. The mesopelagic zone occupies about 60% of the planet's surface and about 20% of the ocean's volume, amounting to a large part of the total biosphere. It hosts a diverse biological community that includes bristlemouths, blobfish, bioluminescent jellyfish, giant squid, and a myriad of other unique organisms adapted to live in a low-light environment. It has long captivated the imagination of scientists, artists and writers; deep sea creatures are prominent in popular culture. Physical conditions The mesopelagic zone includes the region of sharp changes in temperature, salinity and density called the thermocline, halocline, and pycnocline respectively. The temperature variations are large; from over 20 °C (68 °F) at the upper layers to around 4 °C (39 °F) at the boundary with the bathyal zone. The variation in salinity is smaller, typically between 34.5 and 35 psu. The density ranges from 1023 to 1027 g/L of seawater. These changes in temperature, salinity, and density induce stratification which create ocean layers. These different water masses affect gradients and mixing of nutrients and dissolved gasses. This makes this a dynamic zone. The mesopelagic zone has some unique acoustic features. The Sound Fixing and Ranging (SOFAR) channel, where sound travels the slowest due to salinity and temperature variations, is located at the base of the mesopelagic zone at about 600–1,200m. It is a wave-guided zone where sound waves refract within the layer and propagate long distances. The channel got its name during World War II when the US Navy proposed using it as a life saving tool. Shipwreck survivors could drop a small explosive timed to explode in the SOFAR channel and then listening stations could determine the position of the life raft. During the 1950s, the US Navy tried to use this zone to detect Soviet submarines by creating an array of hydrophones called the Sound Surveillance System (SOSUS.) Oceanographers later used this underwater surveillance system to figure out the speed and direction of deep ocean currents by dropping SOFAR floats that could be detected with the SOSUS array. The mesopelagic zone is important for water mass formation, such as mode water. Mode water is a water mass that is typically defined by its vertically mixed properties. It often forms as deep mixed layers at the depth of the thermocline. The mode water in the mesopelagic has residency times on decadal or century scales. The longer overturning times contrast with the daily and shorter scales that a variety of animals move vertically through the zone and sinking of various debris. Biogeochemistry Carbon The mesopelagic zone plays a key role in the ocean's biological pump, which contributes to the oceanic carbon cycle. In the biological pump, organic carbon is produced in the surface euphotic zone where light promotes photosynthesis. A fraction of this production is exported out of the surface mixed layer and into the mesopelagic zone. One pathway for carbon export from the euphotic layer is through sinking of particles, which can be accelerated through repackaging of organic matter in zooplankton fecal pellets, ballasted particles, and aggregates. In the mesopelagic zone, the biological pump is key to carbon cycling, as this zone is largely dominated by remineralization of particulate organic carbon (POC). When a fraction of POC is exported from the euphotic zone, an estimated 90% of that POC is respired in the mesopelagic zone. This is due to the microbial organisms that respire organic matter and remineralize the nutrients, while mesopelagic fish also package organic matter into quick-sinking parcels for deeper export. Another key process occurring in this zone is the diel vertical migration of certain species, which move between the euphotic zone and mesopelagic zone and actively transport particulate organic matter to the deep. In one study in the Equatorial Pacific, myctophids in the mesopelagic zone were estimated to actively transport 15–28% of the passive POC sinking to the deep, while a study near the Canary Islands estimated 53% of vertical carbon flux was due to active transport from a combination of zooplankton and micronekton. When primary productivity is high, the contribution of active transport by vertical migration has been estimated to be comparable to sinking particle export. Particle Packaging and sinking Mean particle sinking rates are 10 to 100 m/day. Sinking rates have been measured in the project VERTIGO (Vertical Transport in the Global Ocean) using settling velocity sediment traps. The variability in sinking rates is due to differences in ballast, water temperature, food web structure and the types of phyto and zooplankton in different areas of the ocean. If the material sinks faster, then it gets respired less by bacteria, transporting more carbon from the surface layer to the deep ocean. Larger fecal pellets sink faster due to lower friction-surface/mass ratio. More viscous waters could slow the sinking rate of particles. Oxygen Dissolved oxygen is a requirement for aerobic respiration, and while the surface ocean is usually oxygen-rich due to atmospheric gas exchange and photosynthesis, the mesopelagic zone is not in direct contact with the atmosphere, due to stratification at the base of the surface mixed layer. Organic matter is exported to the mesopelagic zone from the overlying euphotic layer, while the minimal light in the mesopelagic zone limits photosynthesis. The oxygen consumption due to respiration of most of the sinking organic matter and lack of gas exchange, often creates an oxygen minimum zone (OMZ) in the mesopelagic. The mesopelagic OMZ is particularly severe in the eastern tropical Pacific Ocean and tropical Indian Ocean due to poor ventilation and high rates of organic carbon export to the mesopelagic. Oxygen concentrations in the mesopelagic are occasionally result in suboxic concentrations, making aerobic respiration difficult for organisms. In these anoxic regions, chemosynthesis may occur in which CO2 and reduced compounds such as sulfide or ammonia are taken up to form organic carbon, contributing to the organic carbon reservoir in the mesopelagic. This pathway of carbon fixation has been estimated to be comparable in rate to the contribution by heterotrophic production in this ocean realm. Nitrogen The mesopelagic zone, an area of significant respiration and remineralization of organic particles, is generally nutrient-rich. This is in contrast to the overlying euphotic zone, which is often nutrient-limited. Areas of low oxygen such as OMZ's are a key area of denitrification by prokaryotes, a heterotrophic pathways in which nitrate is converted into nitrogen gas, resulting in a loss to the ocean reservoir of reactive nitrogen. At the suboxic interface that occurs at the edge of the OMZ, nitrite and ammonium can be coupled to produce nitrogen gas through anammox, also removing nitrogen from the biologically available pool. Biology Although some light penetrates the mesopelagic zone, it is insufficient for photosynthesis. The biological community of the mesopelagic zone has adapted to a low-light environment. This is a very efficient ecosystem with many organisms recycling the organic matter sinking from the epipelagic zone resulting in very little organic carbon making it to deeper ocean waters. The general types of life forms found are daytime-visiting herbivores, detritivores feeding on dead organisms and fecal pellets, and carnivores feeding on those detritivores. Many organisms in the mesopelagic zone move up into the epipelagic zone at night, and retreat to the mesopelagic zone during the day, which is known as diel vertical migration. These migrators can therefore avoid visual predators during the day and feed at night, while some of their predators also migrate up at night to follow the prey. There is so much biomass in this migration that sonar operators in World War II would regularly misinterpret the signal returned by this thick layer of plankton as a false sea floor. Estimates of the global biomass of mesopelagic fishes range from 1 gigatonne (Gt) based on net tows to 7–10 Gt based on measurements using active acoustics. Virus and microbial ecology Very little is known about the microbial community of the mesopelagic zone because it is a difficult part of the ocean to study. Recent work using DNA from seawater samples emphasized the importance of viruses and microbes role in recycling organic matter from the surface ocean, known as the microbial loop. These many microbes can get their energy from different metabolic pathways. Some are autotrophs, heterotrophs, and a 2006 study even discovered chemoautotrophs. This chemoautotrophic Archaea crenarchaeon Candidatus can oxidize ammonium as their energy source without oxygen, which could significantly impact the nitrogen and carbon cycles. One study estimates these ammonium-oxidizing bacteria, which are only 5% of the microbial population, can annually capture 1.1 Gt of organic carbon. Microbial biomass and diversity typically decline exponentially with depth in the mesopelagic zone, tracking the general decline of food from above. The community composition varies with depths in the mesopelagic as different organisms are evolved for varying light conditions. Microbial biomass in the mesopelagic is greater at higher latitudes and decreases towards the tropics, which is likely linked to the differing productivity levels in the surface waters. Viruses however are very abundant in the mesopelagic, with around 1010 - 1012 every cubic meter, which is fairly uniform throughout the mesopelagic zone. Zooplankton ecology The mesopelagic zone hosts a diverse zooplankton community. Common zooplankton include copepods, krill, jellyfish, siphonophores, larvaceans, cephalopods, and pteropods. Food is generally scarce in the mesopelagic, so predators have to be efficient in capturing food. Gelatinous organisms are thought to play an important role in the ecology of the mesopelagic and are common predators. Though previously thought to be passive predators just drifting through the water column, jellyfish could be more active predators. One study found that the helmet jellyfish Periphylla periphylla exhibit social behavior and can find each other at depth and form groups. Such behavior was previously attributed to mating, but scientists speculate this could be a feeding strategy to allow a group of jellyfish to hunt together. Mesopelagic zooplankton have unique adaptations for the low light. Bioluminescence is a very common strategy in many zooplankton. This light production is thought to function as a form of communication between conspecifics, prey attraction, prey deterrence, and/or reproduction strategy. Another common adaption are enhanced light organs, or eyes, which is common in krill and shrimp, so they can take advantage of the limited light. Some octopus and krill even have tubular eyes that look upwards in the water column. Most life processes, like growth rates and reproductive rates, are slower in the mesopelagic. Metabolic activity has been shown to decrease with increasing depth and decreasing temperature in colder-water environments. For example, the mesopelagic shrimp-like mysid, Gnathophausia ingens, lives for 6.4 to 8 years, while similar benthic shrimp only live for 2 years. Fish ecology The mesopelagic is home to a significant portion of the world's total fish biomass. Mesopelagic fish are found globally, with exceptions in the Arctic Ocean. A 1980 study puts the mesopelagic fish biomass at about one billion tons. Then a 2008 study estimated the world marine fish biomass at between 0.8 and 2 billion tons. A more recent study concluded mesopelagic fish could have a biomass amounting to 10 billion tons, equivalent to about 100 times the annual catch of traditional fisheries of about 100 million metric tons. However, there is a lot of uncertainty in this biomass estimate. This ocean realm could contain the largest fishery in the world and there is active development for this zone to become a commercial fishery. There are currently thirty families of known mesopelagic fish. One dominant fish in the mesopelagic zone are lanternfish (Myctophidae), which include 245 species distributed among 33 different genera. They have prominent photophores along their ventral side. The Gonostomatidae, or bristlemouth, are also common mesopelagic fish. The bristlemouth could be the Earth's most abundant vertebrate, with numbers in the hundreds of trillions to quadrillions. Mesopelagic fish are difficult to study due to their unique anatomy. Many of these fish have swim bladders to help them control their buoyancy, which makes them hard to sample because those gas-filled chambers typically burst as the fish come up in nets and the fish die. Scientists in California have made progress on mesopelagic fish sampling by developing a submersible chamber that can keep fish alive on the way up to the surface under a controlled atmosphere and pressure. A passive method to estimate mesopelagic fish abundance is by echosounding to locate the 'deep scattering layer' through the backscatter received from these acoustic sounders. A 2015 study suggested that some areas have had a decline in abundance of mesopelagic fish, including off the coast of Southern California, using a long-term study dating back to the 1970s. Cold water species were especially vulnerable to decline. Mesopelagic fish are adapted to a low-light environment. Many fish are black or red, because these colors appear dark due to the limited light penetration at depth. Some fish have rows of photophores, small light-producing organs, on their underside to mimic the surrounding environment. Other fish have mirrored bodies which are angled to reflect the surrounding ocean low-light colors and protect the fish from being seen, while another adaptation is countershading where fish have light colors on the ventral side and dark colors on the dorsal side. Food is often limited and patchy in the mesopelagic, leading to dietary adaptations. Common adaptations fish may have include sensitive eyes and huge jaws for enhanced and opportunistic feeding. Fish are also generally small to reduce the energy requirement for growth and muscle formation. Other feeding adaptations include jaws that can unhinge, elastic throats, and massive, long teeth. Some predators develop bioluminescent lures, like the tasselled anglerfish, which can attract prey, while others respond to pressure or chemical cues instead of relying on vision. Human impacts Pollution Marine debris Marine debris, specifically in the plastic form, have been found in every ocean basin and have a wide range of impacts on the marine world. One of the most critical issues is ingestion of plastic debris, specifically microplastics. Many mesopelagic fish species migrate to the surface waters to feast on their main prey species, zooplankton and phytoplankton, which are mixed with microplastics in the surface waters. Additionally, research has shown that even zooplankton are consuming the microplastics themselves. Mesopelagic fish play a key role in energy dynamics, meaning they provide food to a number of predators including birds, larger fish and marine mammals. The concentration of these plastics has the potential to increase, so more economically important species could become contaminated as well. Concentration of plastic debris in mesopelagic populations can vary depending on geographic location and the concentration of marine debris located there. In 2018, approximately 73% of approximately 200 fish sampled in the North Atlantic had consumed plastic. Bioaccumulation Bioaccumulation (a buildup of a certain substance in the adipose tissue) and biomagnification (the process in which the concentration of the substance grows higher as you rise through the food chain) are growing issues in the mesopelagic zone. Mercury in fish can be traced back to a combination of anthropological factors (such as coal mining) in addition to natural factors. Mercury is a particularly important bioaccumulation contaminant because its concentration in the mesopelagic zone is increasing faster than in surface waters. Inorganic mercury occurs in anthropogenic atmospheric emissions in its gaseous elemental form, which then oxidizes and can be deposited in the ocean. Once there, the oxidized form can be converted to methylmercury, which is its organic form. Research suggests that current levels anthropogenic emissions will not equilibrate between the atmosphere and ocean for a period of decades to centuries, which means we can expect current mercury concentrations in the ocean to keep rising. Mercury is a potent neurotoxin, and poses health risks to the whole food web, beyond the mesopelagic species that consume it. Many of the mesopelagic species, such as myctophids, that make their diel vertical migration to the surface waters, can transfer the neurotoxin when they are consumed by pelagic fish, birds and mammals. Fishing Historically, there have been few examples of efforts to commercialize the mesopelagic zone due to low economic value, technical feasibility and environmental impacts. While the biomass may be abundant, fish species at depth are generally smaller in size and slower to reproduce. Fishing with large trawl nets poses threats to a high percentage of bycatch as well as potential impacts to the carbon cycling processes. Additionally, ships trying to reach productive mesopelagic regions requires fairly long journeys offshore. In 1977, a Soviet fishery opened but closed less than 20 years later due to low commercial profits, while a South African purse seine fishery closed in the mid-1980s due to processing difficulties from the high oil content of fish. As the biomass in the mesopelagic is so abundant, there has been an increased interest to determine whether these populations could be of economic use in sectors other than direct human consumption. For example, it has been suggested that the high abundance of fish in this zone could potentially satisfy a demand for fishmeal and nutraceuticals. With a growing global population, the demand for fishmeal in support of a growing aquaculture industry is high. There is potential for an economically viable harvest. For example, 5 billion tons of mesopelagic biomass could result in the production of circa 1.25 billion tons of food for human consumption. Additionally, the demand for nutraceuticals is also rapidly growing, stemming from the popular human consumption of Omega-3 Fatty Acids in addition to the aquaculture industry that requires a specific marine oil for feed material. Lanternfish are of much interest to the aquaculture market, as they are especially high in fatty acids. Climate Change The mesopelagic region plays an important role in the global carbon cycle, as it is the area where most of the surface organic matter is respired. Mesopelagic species also acquire carbon during their diel vertical migration to feed in surface waters, and they transport that carbon to the deep sea when they die. It is estimated that the mesopelagic cycles between 5 and 12 billion tons of carbon dioxide from the atmosphere per year, and until recently, this estimate was not included in many climate models. It is difficult to quantify the effects of climate change on the mesopelagic zone as a whole, as climate change does not have uniform impacts geographically. Research suggests that in warming waters, as long as there are adequate nutrients and food for fish, then mesopelagic biomass could actually increase due to higher trophic efficiency and increased temperature-driven metabolism. However, because ocean warming will not be uniform throughout the global mesopelagic zone, it is predicted that some areas may actually decrease in fish biomass, while others increase. Water column stratification will also likely increase with ocean warming and climate change. Increased ocean stratification reduces the introduction of nutrients from the deep ocean into the euphotic zone resulting in decreases in both net primary production and sinking particulate matter. Additional research suggests shifts in the geographical range of many species could also occur with warming, with many of them shifting poleward. The combination of these factors could potentially mean that as global ocean basins continue to warm, there could be areas in the mesopelagic that increase in biodiversity and species richness, while declines in other areas, especially moving farther from the equator. Research and Exploration There is a dearth of knowledge about the mesopelagic zone so researchers have begun to develop new technology to explore and sample this area. The Woods Hole Oceanographic Institution (WHOI), NASA, and the Norwegian Institute of Marine Research are all working on projects to gain a better understanding of this zone in the ocean and its influence on the global carbon cycle. Traditional sampling methods like nets have proved to be inadequate because they scare off creatures due to the pressure wave formed by the towed net and the light produced by the bioluminescent species caught in the net. Mesopelagic activity was first investigated by use of sonar because the return bounces off of plankton and fish in the water. However, there are many challenges with acoustic survey methods and previous research has estimated errors in measured amounts of biomass of up to three orders of magnitude. This is due to inaccurate incorporation of depth, species size distribution, and acoustic properties of the species. Norway's Institute of Marine Research has launched a research vessel named Dr. Fridtjof Nansen to investigate mesopelagic activity using sonar with their focus being on the sustainability of fishing operations. To overcome the challenges faced with acoustic sampling, WHOI is developing remote operated vehicles (ROVs) and robots (Deep-See, Mesobot, and Snowclops) that are capable of studying this zone more precisely in a dedicated effort called the Ocean Twilight Zone project that launched in August 2018. Discovery and Detection The deep scattering layer often characterizes the mesopelagic due to the high amount of biomass that exists in the region. Acoustic sound sent into the ocean bounces off particles and organisms in the water column and return a strong signal. The region was initially discovered by American researchers during World War II in 1942 during anti-submarine research with sonar. Sonar at the time could not penetrate below this depth due to the large number of creatures obstructing sound waves. It is uncommon to detect deep scattering layers below 1000m. Until recently, sonar has been the predominant method for studying the mesopelagic. The Malaspina Circumnavigation Expedition was a Spanish-led scientific quest in 2011 to gain a better understanding of the state of the ocean and the diversity in the deep oceans. The data collected, particularly through sonar observations showed that the biomass estimation in the mesopelagic was lower than previously thought. Deep-See WHOI is currently working on a project to characterize and document the pelagic ecosystem. They have developed a device named Deep-See weighing approximately 700 kg, which is designed to be towed behind a research vessel. The Deep-See is capable of reaching depths up to 2000 m and can estimate the amount of biomass and biodiversity in this mesopelagic ecosystem. Deep-See is equipped with cameras, sonars, sensors, water sample collection devices, and a real-time data transmission system. Mesobot WHOI is collaborating with the Monterey Bay Aquarium Research Institute (MBARI), Stanford University, and the University of Texas Rio Grande Valley to develop a small autonomous robot, Mesobot, weighing approximately 75 kg. Mesobot is equipped with high-definition cameras to track and record mesopelagic species on their daily migration over extended periods of time. The robot's thrusters were designed so that they do not disturb the life in the mesopelagic that it is observing. Traditional sample collection devices fail to preserve organisms captured in the mesopelagic due to the large pressure change associated with surfacing. The Mesobot also has a unique sampling mechanism that is capable of keeping the organisms alive during their ascent. The first sea trial of this device is expected to be in 2019. MINIONS Another mesopelagic robot developed by WHOI are the MINIONS. This device descends down the water column and takes images of the amount and size distribution of marine snow at various depths. These tiny particles are a food source for other organisms so it is important to monitor the different levels of marine snow to characterize the carbon cycling processes between the surface ocean and the mesopelagic. SPLAT cam The Harbor Branch Oceanographic Institute has developed the Spatial PLankton Analysis Technique (SPLAT) to identify and map distribution patterns of bioluminescent plankton. The various bioluminescent species produce a unique flash that allows the SPLAT to distinguish each specie's flash characteristic and then map their 3-dimensional distribution patterns. Its intended use was not for investigating the mesopelagic zone, although it is capable of tracking movement patterns of bioluminescent species during their vertical migrations. It would be interesting to apply this mapping technique in the mesopelagic to obtain more information about the diurnal vertical migrations that occur in this zone of the ocean.	Physical sciences	Oceanography	Earth science
565536	https://en.wikipedia.org/wiki/Bathypelagic%20zone	Bathypelagic zone	The bathypelagic zone or bathyal zone (from Greek βαθύς (bathýs), deep) is the part of the open ocean that extends from a depth of below the ocean surface. It lies between the mesopelagic above and the abyssopelagic below. The bathypelagic is also known as the midnight zone because of the lack of sunlight; this feature does not allow for photosynthesis-driven primary production, preventing growth of phytoplankton or aquatic plants. Although larger by volume than the photic zone, human knowledge of the bathypelagic zone remains limited by ability to explore the deep ocean. Physical characteristics The bathypelagic zone is characterized by a nearly constant temperature of approximately and a salinity range of 33-35 g/kg. This region has little to no light because sunlight does not reach this deep in the ocean and bioluminescence is limited. The hydrostatic pressure in this zone ranges from 100-400 atmospheres (atm) due to the increase of 1 atm for every 10 m depth. It is believed that these conditions have been consistent for the past 8000 years. This ocean depth spans from the edge of the continental shelf down to the top of the abyssal zone, and along continental slope depths. The bathymetry of the bathypelagic zone consists of limited areas where the seafloor is in this depth range along the deepest parts of the continental margins, as well as seamounts and mid-ocean ridges. The continental slopes are mostly made up of accumulated sediment, while seamounts and mid-ocean ridges contain large areas of hard substrate that provide habitats for bathypelagic fishes and benthic invertebrates. Although currents at these depths are very slow, the topography of seamounts interrupts the currents and creates eddies that retain plankton in the seamount region, thus increasing fauna nearby as well Hydrothermal vents are also a common feature in some areas of the bathypelagic zone and are primarily formed from the spreading of Earth's tectonic plates at mid-ocean ridges. As the bathypelagic region lacks light, these vents play an important role in global ocean chemical processes, thus supporting unique ecosystems that have adapted to utilize chemicals as energy, via chemoautotrophy, instead of sunlight, to sustain themselves. In addition, hydrothermal vents facilitate precipitation of minerals on the seafloor, making them regions of interest for deep-sea mining. Biogeochemistry Many of the biogeochemical processes in the bathypelagic region are dependent upon the input of organic matter from the overlying epipelagic and mesopelagic zones. This organic material, sometimes called marine snow, sinks in the water column or is transported within downward convected water masses such as the Thermohaline Circulation. Hydrothermal vents also deliver heat and chemicals such as sulfide and methane. These chemicals can be utilized to sustain metabolism by organisms in the region. Our understanding of these biogeochemical processes has historically been limited due to the difficulty and cost of collecting samples from these ocean depths. Other technological challenges, such as measuring microbial activity under the pressure conditions experienced in the bathypelagic zone, have also restricted our knowledge of the region. Although scientific advancements have increased our understanding over the past several decades, many aspects remain a mystery. One of the major areas of current research is focused on understanding carbon remineralization rates in the region. Prior studies have struggled to quantify the rates at which prokaryotes in this region remineralize carbon because previously developed techniques may not be adequate for this region, and indicate remineralization rates much higher than expected. Further work is needed to explore this question, and may require revisions to our understanding of the global carbon cycle. Particulate organic matter Organic material from primary production in the epipelagic zone, and to a far lesser extent, organic inputs from terrestrial sources, make up a majority of the Particulate Organic Matter (POM) in the ocean. POM is delivered to the bathypelagic zone via sinking copepod fecal pellets and dead organisms; these parcels of organic matter fall through the water column and deliver organic carbon, nitrogen, and phosphorus, to organisms that live below the photic zone. These parcels are sometimes referred to as marine snow or ocean dandruff. This is also the dominant delivery mechanism of food to organisms in the bathypelagic zone because there is no sunlight for photosynthesis, with chemoautotrophy playing a more minor role as far as we know. As POM sinks through the water column, it is consumed by organisms which deplete it of nutrients. The size and density of these particles affect their likelihood of reaching organisms in the bathypelagic zone. Smaller parcels of POM often become aggregated together as they fall, which quickens their descent and prohibits their consumption by other organisms, increasing their likelihood of reaching lower depths. The density of these particles may be increased in some regions where minerals associated with some forms of phytoplankton, such as biogenic silica and calcium carbonate "ballast" resulting in more rapid transport to deeper depth. Carbon A majority of organic carbon is produced in the epipelagic zone, with a small portion transported deeper into the ocean interior. This process, known as the biological pump, plays a large role in the sequestration of carbon from the atmosphere into the ocean. Organic carbon is primarily exported to the bathypelagic zone in the form of particulate organic carbon (POC) and dissolved organic carbon (DOC). POC is the largest component of organic carbon delivered to the bathypelagic zone; it primarily takes the form of fecal pellets and dead organisms that sink out of the surface waters and fall toward the ocean floor. Regions with higher primary productivity where particles are able to sink quickly, such as equatorial upwelling zones and the Arabian Sea, have the greatest amount of POC delivery to the bathypelagic zone. The vertical mixing of DOC-rich surface waters is also a process that delivers carbon to the bathypelagic zone, however, it constitutes a substantially smaller portion of overall transport than POC delivery. DOC transport occurs most readily in regions with high rates of ventilation or ocean turnover, such as the interior of gyres or deep water formation sites along the thermohaline circulation. Calcium carbonate dissolution The region in the water column at which calcite dissolution begins to occur rapidly, known as the lysocline, is typically located near the base bathypelagic zone at approximately 3,500 m depth, but varies among ocean basins. The lysocline lies below the saturation depth (the transition to undersaturated conditions with respect to calcium carbonate) and above the carbonate compensation depth (below which there is no calcium carbonate preservation). In a supersaturated environment, the tests of calcite-forming organisms are preserved as they sink toward the sea floor, resulting in sediments with relatively high amounts of CaCO3. However, as depth and pressure increase and temperature decreases, the solubility of calcium carbonate also increases, which results in more dissolution and less net transport to the deeper, underlying seafloor. As a result of this rapid change in dissolution rates, sediments in the bathypelagic region vary widely in CaCO3 content and burial. Ecology The ecology of the bathypelagic ecosystem is constrained by its lack of sunlight and primary producers, with limited production of microbial biomass via autotrophy. The trophic networks in this region rely on particulate organic matter (POM) that sinks from the epipelagic and mesopelagic water, and oxygen inputs from the thermohaline circulation. Despite these limitations, this open-ocean ecosystem is home to microbial organisms, fish, and nekton. Microbial ecology A comprehensive understanding of the inputs driving the microbial ecology in the bathypelagic zone is lacking due to limited observational data, but has been improving with advancements in deep-sea technology. A majority of our knowledge of ocean microbial activity comes from studies of the shallower regions of the ocean because it is easier to access, and it was previously assumed that deeper water did not have suitable physical conditions for diverse microbial communities. The bathypelagic zone receives inputs of organic material and POM from the surface ocean on the order of 1-3.6 Pg C/year. Prokaryote biomass in the bathypelagic is dependent and thus correlated with the amount of sinking POM and organic carbon availability. These essential organic carbon inputs for microbes typically decrease with depth as they are utilized while sinking to the bathypelagic. Microbial production varies over six orders of magnitude based on resource availability in a given area. Prokaryote abundance can range from 0.03-2.3x105 cells ml−1, and have population turnover times that can range from 0.1–30 years. Archaea make up a larger portion of the total prokaryote cell abundance, and different groups have different growth needs, with some archaea groups for example utilizing amino acid groups more readily than others. Some archaea like Crenarchaeota have Crenarchaeota 16S rRNA and archaeal amoA gene abundances correlated to dissolved inorganic carbon (DIC) fixation. The utilization of DIC is thought to be fueled by the oxidation of ammonium and is one form of chemoautotrophy. Based on regional variation and differences in prokaryote abundance, heterotrophic prokaryote production, and particulate organic carbon (POC) inputs to the bathypelagic zone. Research to quantify bacterial-consuming grazers, like heterotrophic eukaryotes, has been limited by difficulties in sampling. Oftentimes organisms do not survive being brought to the surface due to experiencing drastic pressure changes in a short amount of time. Work is underway to quantify cell abundance and biomass, but due to poor survival, it is difficult to get accurate counts. In more recent years there has been an effort to categorize the diversity of the eukaryotic assemblages in the bathypelagic zone using methods to assess the genetic compositions of microbial communities based on supergroups, which is a way to classify organisms that have common ancestry. Some important groups of bacterial grazers include Rhizaria, Alveolata, Fungi, Stramenopiles, Amoebozoa, and Excavata (listed from most to least abundant), with the remaining composition classified as uncertain or other. Viruses influence biogeochemical cycling through the role they play in marine food webs. Their overall abundance can be up to two orders of magnitude lower than the mesopelagic zone, however, there is often high viral abundance found around deep-sea hydrothermal vents. The magnitude of their impacts on biological systems is demonstrated by the varying range of viral-to-prokaryote abundance ratios ranging from 1-223, this indicates that there are the same amount or more viruses than prokaryotes. Fauna Fish ecology Despite the lack of light, vision plays a role in life within the bathypelagic with bioluminescence a trait among both nektonic and planktonic organisms. In contrast to organisms in the water column, benthic organisms in this region tend to have limited to no bioluminescence. The bathypelagic zone contains sharks, squid, octopuses, and many species of fish, including deep-water anglerfish, gulper eel, amphipods, and dragonfish. The fish are characterized by weak muscles, soft skin, and slimy bodies. The adaptations of some of the fish that live there include small eyes and transparent skin. However, this zone is difficult for fish to live in since food is scarce; resulting in species evolving slow metabolic rates in order to conserve energy. Occasionally, large sources of organic matter from decaying organisms, such as whale falls, create a brief burst of activity by attracting organisms from different bathypelagic communities. Diel vertical migration Some bathypelagic species undergo vertical migration, which differs from the diel vertical migration of mesopelagic species in that it is not driven by sunlight. Instead, the migration of bathypelagic organisms is driven by other factors, most of which remain unknown. Some research suggests the movement of species within the overlying pelagic region could prompt individual bathypelagic species to migrate, such as Sthenoteuthis sp., a species of squid. In this particular example, Sthenoteuthis sp. appears to migrate individually over the course of ~4–5 hours towards the surface and then form into groups. While in most regions migration patterns can be driven by predation, in this particular region, the migration patterns are not believed to result solely from predator-prey relations. Instead, these relations are commensalistic, with the species who remain in the bathypelagic benefitting from the POM mixing caused by the upward movement of another species. In addition, the vertical migrating species' timing bathypelagic appears linked to the lunar cycle. However, the exact indicators causing this timing are still unknown. Research and exploration This region is understudied due to a lack of data/observations and difficulty of access (i.e. cost, remote locations, extreme pressure). Historically in oceanography, continental margins were the most sampled and researched due to their relatively easy access. However, more recently locations further offshore and at greater depths, such as ocean ridges and seamounts, are being increasingly studied due to advances in technology and laboratory methods, as well as collaboration with industry. The first discovery of communities subsisting off of the chemical energy in hydrothermal vents was aboard an expedition in 1977 led by Jack Corliss, an oceanographer from Oregon State University. More recent advancements include remotely operated vehicles (ROVs), autonomous underwater vehicles (AUVs), and independent gliders and floats. Specific technologies and research projects SERPENT Project Ocean Twilight Zone (OTZ) Project DEEP SEARCH Project DEEPEND Project AUV Sentry ROV Jason Hybrid ROV Nereus AUV Autosub Long Range Climate change The oceans act as a buffer for anthropogenic climate change due to their ability to take up atmospheric CO2 and absorb heat from the atmosphere. However, the ocean's ability to do so will be negatively affected as atmospheric CO2 concentrations continue to rise and global temperatures continue to warm. This will lead to changes such as deoxygenation, ocean acidification, temperature increase, and carbon sequestration decrease, among other physical and chemical alterations. These perturbations may have significant impacts on the organisms that dwell in the bathypelagic region and the properties that deliver organic carbon to the deep sea. Carbon storage The bathypelagic zone currently acts as a significant reservoir for carbon because of its sheer volume and the century to millennial timescales these waters are isolated from the atmosphere, this ocean zone plays an important role in moderating the effects of anthropogenic climate change. The burial of particulate organic carbon (POC) in the underlying sediments via the biological carbon pump, and the solubility pump of dissolved inorganic carbon (DIC) into the ocean interior via the thermohaline conveyor are key processes for removing excess atmospheric carbon. However, as atmospheric CO2 concentrations and global temperatures continue to rise, the efficiency at which the bathypelagic will store and bury the influx of carbon will most likely decrease. While some regions may experience an increase in POC input, such as Arctic regions where increased periods of minimal sea ice coverage will increase the downward flux of carbon from the surface oceans, overall, there will likely be less carbon sequestered to the bathypelagic region.	Physical sciences	Oceanography	Earth science
566041	https://en.wikipedia.org/wiki/Centaurea%20cyanus	Centaurea cyanus	Centaurea cyanus, commonly known as cornflower or bachelor's button, is an annual flowering plant in the family Asteraceae native to Europe. In the past, it often grew as a weed in cornfields (in the broad sense of "corn", referring to grains, such as wheat, barley, rye, or oats), hence its name. It is now endangered in its native habitat by agricultural intensification, particularly by over-use of herbicides. However, Centaurea cyanus is now also naturalised in many other parts of the world, including North America and parts of Australia through introduction as an ornamental plant in gardens and as a seed contaminant in crop seeds. Description Centaurea cyanus is an annual plant growing to tall, with grey-green branched stems. The leaves are lanceolate and long. The flowers are most commonly an intense blue colour and arranged in flowerheads (capitula) of 1.5–3 cm diameter, with a ring of a few large, spreading ray florets surrounding a central cluster of disc florets. The blue pigment is protocyanin, which in roses is red. Fruits are approx. 3.5 mm long with 2–3 mm-long pappus bristles. It flowers all summer. Genetics Centaurea cyanus is a diploid flower (2n = 24). The genetic diversity within populations is high, although there could be a future decline in diversity due to population fragmentation and intensive agriculture. In general, Centaurea cyanus is a self-incompatible species. However, selfing still occurs occasionally, but results in inbreeding depression. Distribution and habitat Centaurea cyanus is native to temperate Europe, but is widely naturalized outside its native range. It has been present in Britain and Ireland as an archaeophyte (ancient introduction) since the Iron Age. In the United Kingdom, it has declined from 264 sites to just 3 sites in the last 50 years. In reaction to this, the conservation charity Plantlife named it as one of 101 species it would actively work to bring 'back from the brink'. In the County Clare (VC H9) in Ireland, C. cyanus is recorded in arable fields as very rare and almost extinct, while in northeast Ireland, it was abundant before the 1930s. Ecology Weed in arable crops Centaurea cyanus is considered a noxious weed in arable crops, especially cereals and rapeseed. In winter wheat, one plant per m2 can cause a yield loss of up to 30 kg / ha. Centaurea cyanus produces around 800 seed per plant, which are either shed shortly before the harvest of cereals, or they are threshed together with the cereal grains, contributing to the further spread of the species by the harvesting machinery and contaminated seed. The occurrence of Centaurea cyanus strongly decreased during the last decades due to improved seed cleaning, more intensive nitrogen fertilization and herbicide use. However, Centaurea cyanus has become more common in cropland due to an increase in crop rotations dominated by winter cereals and rapeseed and the use of more selective herbicides with a low effectiveness against Centaurea cyanus. In addition, the emergence of resistance against the herbicide class of sulfonylureas has been reported recently. Due to its strong roots, Centaurea cyanus is difficult to control mechanically in spring. Fodder for insects and birds The pollen of Centaurea cyanus is used by several different insect species. Insects of the orders Hymenoptera and Diptera are particularly attracted by the flower. As Centaurea cyanus is a self-incompatible species, it needs external pollination. The nectar of Centaurea cyanus is very sweet with a sugar content of 34%. Due to its high sugar production of up to 0.2 mg sugar per day and flower, the species is highly appreciated by beekeepers. The seeds of Centaurea cyanus are one of the favourite foods of the European goldfinch. Control of insect pests Centaurea cyanus was found to produce volatiles attracting Microplitis mediator, which is a major parasitoid of the cabbage moth (Mamestra brassicae), which is the most important pest of cabbage (Brassica oleracea) in central Europe. Planting Centaurea cyanus in cabbage fields as a companion plant was thus suggested as an alternative to the widespread use of insecticides to control Mamestra brassicae. Field experiments showed that planting Centaurea cyanus in cabbage fields at a density of 1 plant / m2 can result in a significant increase in parasitation of Mamestra brassicae larvae, predation of Mamestra brassicae eggs (e.g. by carabid beetles or spiders) and ultimately cabbage yield. Cultivation Several cultivars of Centaurea cyanus with varying pastel colours, including pink and purple, have been selected for ornamental purposes. The species is also grown for the cut flower industry in Canada for use by florists. Doubled blue cultivars (such as 'Blue Boy' or 'Blue Diadem') are most commonly used for this purpose, but white, pink, lavender and black (actually a very dark maroon) cultivars are also used, albeit to a lesser extent. There are varieties with blue, white, purple, pink or even black petals. Breeding goals As for all ornamental plants, important goals of Centaurea cyanus breeding include the induction of phenotypic variation (e.g. in flower coloration, size and shape, foliage characteristics or plant height), higher flower yield, resistance to pests and diseases as well as tolerance to abiotic stress (e.g., extreme temperatures, drought or salinity). Soil and climate requirements Centaurea cyanus requires full sun and neutral (pH 6.6–7.5) to mildly alkaline (pH 7.6–7.8), moist and well-drained soil. However, Centaurea cyanus is quite tolerant to drought once established. Sowing For summer-blooming plants, sowing should be executed in late spring. In moderate climates, however, it is also possible to sow Centaurea cyanus in early fall. In this case, plants will already start to flower in the following spring. Recommended spacing between plants is approx. 20 to 30 cm. Centaurea cyanus can germinate from up to 10 cm depth, but the best result is obtained at 1 cm sowing depth. Germination occurs quickly after sowing. Fertilization and cultural practices High phosphorus fertilization in mid-summer will increase flower production. Mulching is recommended to prevent drying out of the soil and exposure of the root system to the sun. Pests and diseases In general, Centaurea cyanus is not very susceptible to pests and plant diseases. However, it may be affected by stem rot and stem rust if grown too tightly or by powdery mildew. Furthermore, aphids and leafhoppers can cause relevant damage to Centaurea cyanus. Seed harvesting Seeds are harvested either by hand or, in an agricultural setting, with a seed harvesting machine. On average there are 97,000 seeds in a pound of cornflower seeds. Hand collecting can be time-consuming and yields are rather low. A seed harvesting machine is more efficient than collecting the seeds by hand, but it is costly. The main principle of such a machine is that it brushes the ripe seeds off the plant and creates a cross flow fan action that generates sufficient air velocity to hold and gather the seeds into the seed bunker. Pruning Deadheading will encourage the plant to produce more blooms. Cornflowers are often used for ornamental purposes and by cutting them, up to their third leaves, they will produce more blooms and grow a bigger stem. Uses Culinary The flowers of Centaurea cyanus can be eaten raw, dried or cooked. Dried petals are used in foods, including in spices. Their main purpose is to add colour to food. There are cheeses or oils that contain raw petals. Petals can also be added to salads, drinks, and desserts for garnishing purposes in raw or dried form. Dried petals are also used in teas and other beverages. Blue cornflower petals are sometimes one of the ingredients in Lady Grey tea. Medicine Centaurea cyanus contains a wide range of pharmacologically active compounds, such as flavonoids, anthocyanins and aromatic acids. Especially the flower head finds application in herbal medicine, but leaves and seeds are also used for pharmacological purposes, albeit to a lesser extent. In particular, extracts from the flower heads have anti-inflammatory properties used in the treatment of minor ocular inflammations. Antioxidant properties are high due to ascorbic acid and phenolic compounds. Furthermore, extracts of the flower head and vegetative parts of the plant were shown to have gastroprotective effects due to their content of quercetin, apigenin and caffeic acid derivates. Pigment The blue color of Centaurea cyanus is due to protocyanin, an anthocyanin pigment that is also found in roses. Different anthocyanins derived from Centaurea cyanus are used as natural additives in food products, such as yoghurts. Phytoremediation Centaurea cyanus has been evaluated for phytoremediation of soils contaminated with lead. Inoculation of the contaminated soil with Glomus spp. (fungus) and Pseudomonas spp. (bacterium) would significantly enhance the biomass production and lead uptake of Centaurea cyanus. In culture In folklore, cornflowers were worn by young men in love; if the flower faded too quickly, it was taken as a sign that the man's love was not returned. The blue cornflower was one of the national symbols of Germany. This is partly due to the story that when Queen Louise of Prussia was fleeing Berlin and pursued by Napoleon's forces, she hid her children in a field of cornflowers and kept them quiet by weaving wreaths for them from the flowers. The flower thus became identified with Prussia, not least because it was the same color as the Prussian military uniform. After the unification of Germany in 1871, it went on to become a symbol of the country as a whole. For this reason, in Austria the blue cornflower is a political symbol for pan-German and rightist ideas. It was worn as a secret symbol identifying members of the then-illegal NSDAP in Austria in the 1930s. Members of the Freedom Party wore it at the openings of the Austrian parliament since 2006. After the last general election 2017 they replaced it with the edelweiss. It was also the favourite flower of Louise's son Kaiser Wilhelm I. Because of its ties to royalty, authors such as Theodor Fontane have used it symbolically, often sarcastically, to comment on the social and political climate of the time. The cornflower is also often seen as an inspiration for the German Romantic symbol of the Blue Flower. Due to its traditional association with Germany, the cornflower has been made the official symbol of the annual German-American Steuben Parade. The blue cornflower has been the national flower of Estonia since 1969 and symbolizes daily bread to Estonians. It is also the symbol of the Estonian Conservative People's Party. It is also the symbol of the Finnish National Coalition Party, and the Liberal People's Party of Sweden, where it has since the dawn of the 20th century been a symbol for social liberalism. It is the official flower of the Swedish province of Östergötland and the school flower of Winchester College and also of Dulwich College, where it is said to have been the favourite flower of the founder, Edward Alleyn. In France the is the symbol of the 11 November 1918 armistice and, as such, a common symbol for veterans (especially the now defunct poilus of World War I), similar to the Remembrance poppies worn in the United Kingdom and in Canada. The cornflower is also the symbol for motor neurone disease and amyotrophic lateral sclerosis. Cornflowers are sometimes worn by Old Harrovians, former pupils of the British Harrow School. A blue cornflower was used by Corning Glass Works for the initial release of Corning Ware Pyroceram cookware. Its popularity in the United States, Canada, United Kingdom and Australia was so high that it became the symbol of Corning Glass Works. In paintings	Biology and health sciences	Asterales	Plants
14355756	https://en.wikipedia.org/wiki/Metal%E2%80%93semiconductor%20junction	Metal–semiconductor junction	In solid-state physics, a metal–semiconductor (M–S) junction is a type of electrical junction in which a metal comes in close contact with a semiconductor material. It is the oldest type of practical semiconductor device. M–S junctions can either be rectifying or non-rectifying. The rectifying metal–semiconductor junction forms a Schottky barrier, making a device known as a Schottky diode, while the non-rectifying junction is called an ohmic contact. (In contrast, a rectifying semiconductor–semiconductor junction, the most common semiconductor device today, is known as a p–n junction.) Metal–semiconductor junctions are crucial to the operation of all semiconductor devices. Usually an ohmic contact is desired, so that electrical charge can be conducted easily between the active region of a transistor and the external circuitry. Occasionally however a Schottky barrier is useful, as in Schottky diodes, Schottky transistors, and metal–semiconductor field effect transistors. The critical parameter: Schottky barrier height Whether a given metal-semiconductor junction is an ohmic contact or a Schottky barrier depends on the Schottky barrier height, ΦB, of the junction. For a sufficiently large Schottky barrier height, that is, ΦB is significantly higher than the thermal energy kT, the semiconductor is depleted near the metal and behaves as a Schottky barrier. For lower Schottky barrier heights, the semiconductor is not depleted and instead forms an ohmic contact to the metal. The Schottky barrier height is defined differently for n-type and p-type semiconductors (being measured from the conduction band edge and valence band edge, respectively). The alignment of the semiconductor's bands near the junction is typically independent of the semiconductor's doping level, so the n-type and p-type Schottky barrier heights are ideally related to each other by: where Eg is the semiconductor's band gap. In practice, the Schottky barrier height is not precisely constant across the interface, and varies over the interfacial surface. Schottky–Mott rule and Fermi level pinning The Schottky–Mott rule of Schottky barrier formation, named for Walter H. Schottky and Nevill Mott, predicts the Schottky barrier height based on the vacuum work function of the metal relative to the vacuum electron affinity (or vacuum ionization energy) of the semiconductor: This model is derived based on the thought experiment of bringing together the two materials in vacuum, and is closely related in logic to Anderson's rule for semiconductor-semiconductor junctions. Different semiconductors respect the Schottky–Mott rule to varying degrees. Although the Schottky–Mott model correctly predicted the existence of band bending in the semiconductor, it was found experimentally that it would give grossly incorrect predictions for the height of the Schottky barrier. A phenomenon referred to as "Fermi level pinning" caused some point of the band gap, at which finite DOS exists, to be locked (pinned) to the Fermi level. This made the Schottky barrier height almost completely insensitive to the metal's work function: where Ebandgap is the size of band gap in the semiconductor. In fact, empirically, it is found that neither of the above extremes is quite correct. The choice of metal does have some effect, and there appears to be a weak correlation between the metal work function and the barrier height, however the influence of the work function is only a fraction of that predicted by the Schottky-Mott rule. It was noted in 1947 by John Bardeen that the Fermi level pinning phenomenon would naturally arise if there were chargeable states in the semiconductor right at the interface, with energies inside the semiconductor's gap. These would either be induced during the direct chemical bonding of the metal and semiconductor (metal-induced gap states) or be already present in the semiconductor–vacuum surface (surface states). These highly dense surface states would be able to absorb a large quantity of charge donated from the metal, effectively shielding the semiconductor from the details of the metal. As a result, the semiconductor's bands would necessarily align to a location relative to the surface states which are in turn pinned to the Fermi level (due to their high density), all without influence from the metal. The Fermi level pinning effect is strong in many commercially important semiconductors (Si, Ge, GaAs), and thus can be problematic for the design of semiconductor devices. For example, nearly all metals form a significant Schottky barrier to n-type germanium and an ohmic contact to p-type germanium, since the valence band edge is strongly pinned to the metal's Fermi level. The solution to this inflexibility requires additional processing steps such as adding an intermediate insulating layer to unpin the bands. (In the case of germanium, germanium nitride has been used) History The rectification property of metal–semiconductor contacts was discovered by Ferdinand Braun in 1874 using mercury metal contacted with copper sulfide and iron sulfide semiconductors. Sir Jagadish Chandra Bose applied for a US patent for a metal-semiconductor diode in 1901. This patent was awarded in 1904. G.W. Pickard received a patent in 1906 on a point-contact rectifier using silicon. In 1907, George W. Pierce published a paper in Physical Review showing rectification properties of diodes made by sputtering many metals on many semiconductors. The use of the metal–semiconductor diode rectifier was proposed by Lilienfeld in 1926 in the first of his three transistor patents as the gate of the metal–semiconductor field effect transistors. The theory of the field-effect transistor using a metal/semiconductor gate was advanced by William Shockley in 1939. The earliest metal–semiconductor diodes in electronics application occurred around 1900, when the cat's whisker rectifiers were used in receivers. They consisted of pointed tungsten wire (in the shape of a cat's whisker) whose tip or point was pressed against the surface of a galena (lead sulfide) crystal. The first large area rectifier appeared around 1926 which consisted of a copper(I) oxide semiconductor thermally grown on a copper substrate. Subsequently, selenium films were evaporated onto large metal substrates to form the rectifying diodes. These selenium rectifiers were used (and are still used) to convert alternating current to direct current in electrical power applications. During 1925–1940, diodes consisting of a pointed tungsten metal wire in contact with a silicon crystal base, were fabricated in laboratories to detect microwaves in the UHF range. A World War II program to manufacture high-purity silicon as the crystal base for the point-contact rectifier was suggested by Frederick Seitz in 1942 and successfully undertaken by the Experimental Station of the E. I du Pont de Nemours Company. The first theory that predicted the correct direction of rectification of the metal–semiconductor junction was given by Nevill Mott in 1939. He found the solution for both the diffusion and drift currents of the majority carriers through the semiconductor surface space charge layer which has been known since about 1948 as the Mott barrier. Walter H. Schottky and Spenke extended Mott's theory by including a donor ion whose density is spatially constant through the semiconductor surface layer. This changed the constant electric field assumed by Mott to a linearly decaying electric field. This semiconductor space-charge layer under the metal is known as the Schottky barrier. A similar theory was also proposed by Davydov in 1939. Although it gives the correct direction of rectification, it has also been proven that the Mott theory and its Schottky-Davydov extension gives the wrong current limiting mechanism and wrong current-voltage formulae in silicon metal/semiconductor diode rectifiers. The correct theory was developed by Hans Bethe and reported by him in a M.I.T. Radiation Laboratory Report dated November 23, 1942. In Bethe's theory, the current is limited by thermionic emission of electrons over the metal–semiconductor potential barrier. Thus, the appropriate name for the metal–semiconductor diode should be the Bethe diode, instead of the Schottky diode, since the Schottky theory does not predict the modern metal–semiconductor diode characteristics correctly. If a metal-semiconductor junction is formed by placing a droplet of mercury, as Braun did, onto a semiconductor, e.g. silicon, to form a Schottky barrier in a Schottky diode electrical setup – electrowetting can be observed, where the droplet spreads out with increasing voltage. Depending on the doping type and density in the semiconductor, the droplet spreading depends on the magnitude and sign of the voltage applied to the mercury droplet. This effect has been termed ‘Schottky electrowetting’, effectively linking electrowetting and semiconductor effects. Between 1953-1958, Fuller and Ditzenberger's work on the diffusion of impurities into silicon. In 1956 Miller and Savage studied the diffusion of aluminium in crystal silicon. The first silicon oxide gate transistor were invented by Frosch and Derick in 1957 at Bell Labs. In 1956, Richard Baker described some discrete diode clamp circuits to keep transistors from saturating. The circuits are now known as Baker clamps. One of those clamp circuits used a single germanium diode to clamp a silicon transistor in a circuit configuration that is the same as the Schottky transistor. The circuit relied on the germanium diode having a lower forward voltage drop than a silicon diode would have. The Schottky diode, also known as the Schottky-barrier diode, was theorized for years, but was first practically realized as a result of the work of Atalla and Kahng during 19601961. They published their results in 1962 and called their device the "hot electron" triode structure with semiconductor-metal emitter. It was one of the first metal-base transistors. Atalla continued research on Schottky diodes with Robert J. Archer at HP Associates. They developed high vacuum metal film deposition technology, and fabricated stable evaporated/sputtered contacts, publishing their results in January 1963. Their work was a breakthrough in metal–semiconductor junction and Schottky barrier research, as it overcame most of the fabrication problems inherent in point-contact diodes and made it possible to build practical Schottky diodes. In 1967, Robert Kerwin, Donald Klein and John Sarace at Bell Labs, patented a method to replaced the aluminum gate with a polycrystalline layer of silicon.	Physical sciences	Electrical circuits	Physics
13201685	https://en.wikipedia.org/wiki/Smart%20grid	Smart grid	The smart grid is an enhancement of the 20th century electrical grid, using two-way communications and distributed so-called intelligent devices. Two-way flows of electricity and information could improve the delivery network. Research is mainly focused on three systems of a smart grid – the infrastructure system, the management system, and the protection system. Electronic power conditioning and control of the production and distribution of electricity are important aspects of the smart grid. The smart grid represents the full suite of current and proposed responses to the challenges of electricity supply. Numerous contributions to the overall improvement of energy infrastructure efficiency are anticipated from the deployment of smart grid technology, in particular including demand-side management. The improved flexibility of the smart grid permits greater penetration of highly variable renewable energy sources such as solar power and wind power, even without the addition of energy storage. Smart grids could also monitor/control residential devices that are noncritical during periods of peak power consumption, and return their function during nonpeak hours. A smart grid includes a variety of operation and energy measures: Advanced metering infrastructure (of which smart meters are a generic name for any utility side device even if it is more capable e.g. a fiber optic router) Smart distribution boards and circuit breakers integrated with home control and demand response (behind the meter from a utility perspective) Load control switches and smart appliances, often financed by efficiency gains on municipal programs (e.g. PACE financing) Renewable energy resources, including the capacity to charge parked (electric vehicle) batteries or larger arrays of batteries recycled from these, or other energy storage. Energy efficient resources Electric surplus distribution by power lines and auto-smart switch Sufficient utility grade fiber broadband to connect and monitor the above, with wireless as a backup. Sufficient spare if "dark" capacity to ensure failover, often leased for revenue. Concerns with smart grid technology mostly focus on smart meters, items enabled by them, and general security issues. Roll-out of smart grid technology also implies a fundamental re-engineering of the electricity services industry, although typical usage of the term is focused on the technical infrastructure. Smart grid policy is organized in Europe as Smart Grid European Technology Platform. Policy in the United States is described in Title 42 of the United States Code. Background Historical development of the electricity grid The first alternating current power grid system was installed in 1886 in Great Barrington, Massachusetts. At that time, the grid was a centralized unidirectional system of electric power transmission, electricity distribution, and demand-driven control. In the 20th century, local grids grew over time and were eventually interconnected for economic and reliability reasons. By the 1960s, the electric grids of developed countries had become very large, mature, and highly interconnected, with thousands of 'central' generation power stations delivering power to major load centres via high capacity power lines which were then branched and divided to provide power to smaller industrial and domestic users over the entire supply area. The topology of the 1960s grid was a result of the strong economies of scale: large coal-, gas- and oil-fired power stations in the 1 GW (1000 MW) to 3 GW scale are still found to be cost-effective, due to efficiency-boosting features that can be cost-effective only when the stations become very large. Power stations were located strategically to be close to fossil fuel reserves (either the mines or wells themselves or else close to rail, road, or port supply lines). Siting of hydroelectric dams in mountain areas also strongly influenced the structure of the emerging grid. Nuclear power plants were sited for the availability of cooling water. Finally, fossil fuel-fired power stations were initially very polluting and were sited as far as economically possible from population centres once electricity distribution networks permitted it. By the late 1960s, the electricity grid reached the overwhelming majority of the population of developed countries, with only outlying regional areas remaining 'off-grid'. Metering of electricity consumption was necessary on a per-user basis in order to allow appropriate billing according to the (highly variable) level of consumption of different users. Because of limited data collection and processing capability during the period of growth of the grid, fixed-tariff arrangements were commonly put in place, as well as dual-tariff arrangements where night-time power was charged at a lower rate than daytime power. The motivation for dual-tariff arrangements was the lower night-time demand. Dual tariffs made possible the use of low-cost night-time electrical power in applications such as the maintaining of 'heat banks' which served to 'smooth out' the daily demand, and reduce the number of turbines that needed to be turned off overnight, thereby improving the utilisation and profitability of the generation and transmission facilities. The metering capabilities of the 1960s grid meant technological limitations on the degree to which price signals could be propagated through the system. From the 1970s to the 1990s, growing demand led to increasing numbers of power stations. In some areas, the supply of electricity, especially at peak times, could not keep up with this demand, resulting in poor power quality including blackouts, power cuts, and brownouts. Increasingly, electricity was depended on for industry, heating, communication, lighting, and entertainment, and consumers demanded ever-higher levels of reliability. Towards the end of the 20th century, electricity demand patterns were established: domestic heating and air-conditioning led to daily peaks in demand that were met by an array of 'peaking power generators' that would only be turned on for short periods each day. The relatively low utilisation of these peaking generators (commonly, gas turbines were used due to their relatively lower capital cost and faster start-up times), together with the necessary redundancy in the electricity grid, resulting in high costs to the electricity companies, which were passed on in the form of increased tariffs. In the 21st century, some developing countries like China, India, and Brazil were seen as pioneers of smart grid deployment. Modernization opportunities Since the early 21st century, opportunities to take advantage of improvements in electronic communication technology to resolve the limitations and costs of the electrical grid have become apparent. Technological limitations on metering no longer force peak power prices to be averaged out and passed on to all consumers equally. In parallel, growing concerns over environmental damage from fossil-fired power stations have led to a desire to use large amounts of renewable energy. Dominant forms such as wind power and solar power are highly variable, and so the need for more sophisticated control systems became apparent, to facilitate the connection of sources to the otherwise highly controllable grid. Power from photovoltaic cells (and to lesser extent wind turbines) has also, significantly, called into question the imperative for large, centralised power stations. The rapidly falling costs point to a major change from the centralised grid topology to one that is highly distributed, with power being both generated and consumed right at the limits of the grid. Finally, growing concern over terrorist attacks in some countries has led to calls for a more robust energy grid that is less dependent on centralised power stations that were perceived to be potential attack targets. Definition of "smart grid" United States The first official definition of Smart Grid was provided by the Energy Independence and Security Act of 2007 (EISA-2007), which was approved by the US Congress in January 2007, and signed to law by President George W. Bush in December 2007. Title XIII of this bill provides a description, with ten characteristics, that can be considered a definition for Smart Grid, as follows:"It is the policy of the United States to support the modernization of the Nation's electricity transmission and distribution system to maintain a reliable and secure electricity infrastructure that can meet future demand growth and to achieve each of the following, which together characterize a Smart Grid: (1) Increased use of digital information and controls technology to improve reliability, security, and efficiency of the electric grid. (2) Dynamic optimization of grid operations and resources, with full cyber-security. (3) Deployment and integration of distributed resources and generation, including renewable resources. (4) Development and incorporation of demand response, demand-side resources, and energy-efficiency resources. (5) Deployment of 'smart' technologies (real-time, automated, interactive technologies that optimize the physical operation of appliances and consumer devices) for metering, communications concerning grid operations and status, and distribution automation. (6) Integration of 'smart' appliances and consumer devices. (7) Deployment and integration of advanced electricity storage and peak-shaving technologies, including plug-in electric and hybrid electric vehicles, and thermal storage air conditioning. (8) Provision to consumers of timely information and control options. (9) Development of standards for communication and interoperability of appliances and equipment connected to the electric grid, including the infrastructure serving the grid. (10) Identification and lowering of unreasonable or unnecessary barriers to adoption of smart grid technologies, practices, and services." European Union The European Union Commission Task Force for Smart Grids also provides smart grid definition as: "A Smart Grid is an electricity network that can cost efficiently integrate the behaviour and actions of all users connected to it – generators, consumers and those that do both – in order to ensure economically efficient, sustainable power system with low losses and high levels of quality and security of supply and safety. A smart grid employs innovative products and services together with intelligent monitoring, control, communication, and self-healing technologies in order to: Better facilitate the connection and operation of generators of all sizes and technologies. Allow consumers to play a part in optimising the operation of the system. Provide consumers with greater information and options for how they use their supply. Significantly reduce the environmental impact of the whole electricity supply system. Maintain or even improve the existing high levels of system reliability, quality and security of supply. Maintain and improve the existing services efficiently." That definition was used in the European Commission Communication (2011) 202. A common element to most definitions is the application of digital processing and communications to the power grid, making data flow and information management central to the smart grid. Various capabilities result from the deeply integrated use of digital technology with power grids. Integration of the new grid information is one of the key issues in the design of smart grids. Electric utilities now find themselves making three classes of transformations: improvement of infrastructure, called the strong grid in China; addition of the digital layer, which is the essence of the smart grid; and business process transformation, necessary to capitalize on the investments in smart technology. Much of the work that has been going on in electric grid modernization, especially substation and distribution automation, is now included in the general concept of the smart grid. Early technological innovations Smart grid technologies emerged from earlier attempts at using electronic control, metering, and monitoring. In the 1980s, automatic meter reading was used for monitoring loads from large customers and evolved into the Advanced Metering Infrastructure of the 1990s, whose meters could store how electricity was used at different times of the day. Smart meters add continuous communications so that monitoring can be done in real-time, and can be used as a gateway to demand response-aware devices and "smart sockets" in the home. Early forms of such demand side management technologies were dynamic demand aware devices that passively sensed the load on the grid by monitoring changes in the power supply frequency. Devices such as industrial and domestic air conditioners, refrigerators, and heaters adjusted their duty cycle to avoid activation during times the grid was suffering a peak condition. Beginning in 2000, Italy's Telegestore Project was the first to network large numbers (27 million) of homes using smart meters connected via low bandwidth power line communication. Some experiments used the term broadband over power lines (BPL), while others used wireless technologies such as mesh networking promoted for more reliable connections to disparate devices in the home as well as supporting metering of other utilities such as gas and water. Monitoring and synchronization of wide-area networks were revolutionized in the early 1990s when the Bonneville Power Administration expanded its smart grid research with prototype sensors that are capable of very rapid analysis of anomalies in electricity quality over very large geographic areas. The culmination of this work was the first operational Wide Area Measurement System (WAMS) in 2000. Other countries are rapidly integrating this technology — China started having a comprehensive national WAMS when the past 5-year economic plan was completed in 2012. The earliest deployments of smart grids include the Italian system Telegestore (2005), the mesh network of Austin, Texas (since 2003), and the smart grid in Boulder, Colorado (2008). See below. Features A smart grid would allow the power industry to observe and control parts of the system at higher resolution in time and space. One of the purposes of the smart grid is real time information exchange to make operation as efficient as possible. It would allow management of the grid on all time scales from high-frequency switching devices on a microsecond scale, to wind and solar output variations on a minute scale, to the future effects of the carbon emissions generated by power production on a decade scale. The smart grid represents the full suite of current and proposed responses to the challenges of electricity supply. Because of the diverse range of factors, there are numerous competing taxonomies and no agreement on a universal definition. Nevertheless, one possible categorization is given here. Reliability The smart grid makes use of technologies such as state estimation, that improve fault detection and allow self-healing of the network without the intervention of technicians. This will ensure a more reliable supply of electricity and reduce vulnerability to natural disasters or attacks. Although multiple routes are touted as a feature of the smart grid, the old grid also featured multiple routes. Initial power lines in the grid were built using a radial model, later connectivity was guaranteed via multiple routes, referred to as a network structure. However, this created a new problem: if the current flow or related effects across the network exceed the limits of any particular network element, it could fail, and the current would be shunted to other network elements, which eventually may fail also, causing a domino effect. See power outage. A technique to prevent this is load shedding by rolling blackout or voltage reduction (brownout). Flexibility in network topology Next-generation transmission and distribution infrastructure will be better able to handle possible bidirectional energy flows, allowing for distributed generation such as from photovoltaic panels on building roofs, but also charging to/from the batteries of electric cars, wind turbines, pumped hydroelectric power, the use of fuel cells, and other sources. Classic grids were designed for a one-way flow of electricity, but if a local sub-network generates more power than it is consuming, the reverse flow can raise safety and reliability issues. A smart grid aims to manage these situations. Efficiency Numerous contributions to the overall improvement of the efficiency of energy infrastructure are anticipated from the deployment of smart grid technology, in particular including demand-side management, for example turning off air conditioners during short-term spikes in electricity price, reducing the voltage when possible on distribution lines through Voltage/VAR Optimization (VVO), eliminating truck-rolls for meter reading, and reducing truck-rolls by improved outage management using data from Advanced Metering Infrastructure systems. The overall effect is less redundancy in transmission and distribution lines, and greater utilization of generators, leading to lower power prices. Load adjustment/Load balancing The total load connected to the power grid can vary significantly over time. Although the total load is the sum of many individual choices of the clients, the overall load is not necessarily stable or slow varying. For example, if a popular television program starts, millions of televisions will start to draw current instantly. Traditionally, to respond to a rapid increase in power consumption, faster than the start-up time of a large generator, some spare generators are put on a dissipative standby mode. A smart grid may warn all individual television sets, or another larger customer, to reduce the load temporarily (to allow time to start up a larger generator) or continuously (in the case of limited resources). Using mathematical prediction algorithms it is possible to predict how many standby generators need to be used, to reach a certain failure rate. In the traditional grid, the failure rate can only be reduced at the cost of more standby generators. In a smart grid, the load reduction by even a small portion of the clients may eliminate the problem. Peak curtailment/leveling and time of use pricing To reduce demand during the high-cost peak usage periods, communications and metering technologies inform smart devices in the home and business when energy demand is high and track how much electricity is used and when it is used. It also gives utility companies the ability to reduce consumption by communicating to devices directly in order to prevent system overloads. Examples would be a utility reducing the usage of a group of electric vehicle charging stations or shifting temperature set points of air conditioners in a city. To motivate them to cut back use and perform what is called peak curtailment or peak leveling, prices of electricity are increased during high demand periods and decreased during low demand periods. It is thought that consumers and businesses will tend to consume less during high-demand periods if it is possible for consumers and consumer devices to be aware of the high price premium for using electricity at peak periods. This could mean making trade-offs such as cycling on/off air conditioners or running dishwashers at 9 pm instead of 5 pm. When businesses and consumers see a direct economic benefit of using energy at off-peak times, the theory is that they will include the energy cost of operation into their consumer device and building construction decisions and hence become more energy efficient. Sustainability The improved flexibility of the smart grid permits greater penetration of highly variable renewable energy sources such as solar power and wind power, even without the addition of energy storage. Current network infrastructure is not built to allow for many distributed feed-in points, and typically even if some feed-in is allowed at the local (distribution) level, the transmission-level infrastructure cannot accommodate it. Rapid fluctuations in distributed generation, such as due to cloudy or gusty weather, present significant challenges to power engineers who need to ensure stable power levels through varying the output of the more controllable generators such as gas turbines and hydroelectric generators. Smart grid technology is a necessary condition for very large amounts of renewable electricity on the grid for this reason. There is also support for vehicle-to-grid. Market-enabling The smart grid allows for systematic communication between suppliers (their energy price) and consumers (their willingness-to-pay), and permits both the suppliers and the consumers to be more flexible and sophisticated in their operational strategies. Only the critical loads will need to pay the peak energy prices, and consumers will be able to be more strategic in when they use energy. Generators with greater flexibility will be able to sell energy strategically for maximum profit, whereas inflexible generators such as base-load steam turbines and wind turbines will receive a varying tariff based on the level of demand and the status of the other generators currently operating. The overall effect is a signal that awards energy efficiency, and energy consumption that is sensitive to the time-varying limitations of the supply. At the domestic level, appliances with a degree of energy storage or thermal mass (such as refrigerators, heat banks, and heat pumps) will be well placed to 'play' the market and seek to minimise energy cost by adapting demand to the lower-cost energy support periods. This is an extension of the dual-tariff energy pricing mentioned above. Demand response support Demand response support allows generators and loads to interact in an automated fashion in real-time, coordinating demand to flatten spikes. Eliminating the fraction of demand that occurs in these spikes eliminates the cost of adding reserve generators, cuts wear and tear and extends the life of equipment, and allows users to cut their energy bills by telling low priority devices to use energy only when it is cheapest. Currently, power grid systems have varying degrees of communication within control systems for their high-value assets, such as in generating plants, transmission lines, substations, and major energy users. In general, information flows one way, from the users and the loads they control back to the utilities. The utilities attempt to meet the demand and succeed or fail to varying degrees (brownouts, rolling blackout, uncontrolled blackout). The total amount of power demanded by the users can have a very wide probability distribution which requires spare generating plants in standby mode to respond to the rapidly changing power usage. This one-way flow of information is expensive; the last 10% of generating capacity may be required as little as 1% of the time, and brownouts and outages can be costly to consumers. Demand response can be provided by commercial, residential loads, and industrial loads. For example, Alcoa's Warrick Operation is participating in MISO as a qualified Demand Response Resource, and the Trimet Aluminium uses its smelter as a short-term mega-battery. Latency of the data flow is a major concern, with some early smart meter architectures allowing actually as long as 24 hours delay in receiving the data, preventing any possible reaction by either supplying or demanding devices. Technology The bulk of smart grid technologies are already used in other applications such as manufacturing and telecommunications and are being adapted for use in grid operations. Integrated communications: Areas for improvement include: substation automation, demand response, distribution automation, supervisory control, and data acquisition (SCADA), energy management systems, wireless mesh networks and other technologies, power-line carrier communications, and fiber-optics. Integrated communications will allow for real-time control, information, and data exchange to optimize system reliability, asset utilization, and security. Sensing and measurement: core duties are evaluating congestion and grid stability, monitoring equipment health, energy theft prevention, and control strategies support. Technologies include advanced microprocessor meters (smart meter) and meter reading equipment, wide-area monitoring systems, (typically based on online readings by Distributed temperature sensing combined with Real time thermal rating (RTTR) systems), electromagnetic signature measurement/analysis, time-of-use, and real-time pricing tools, advanced switches and cables, backscatter radio technology, and Digital protective relays. Smart meters. Phasor measurement units. Many in the power systems engineering community believe that the Northeast blackout of 2003 could have been contained to a much smaller area if a wide area phasor measurement network had been in place. Distributed power flow control: power flow control devices clamp onto existing transmission lines to control the flow of power within. Transmission lines enabled with such devices support greater use of renewable energy by providing more consistent, real-time control over how that energy is routed within the grid. This technology enables the grid to more effectively store intermittent energy from renewables for later use. Smart power generation using advanced components: smart power generation is a concept of matching electricity generation with demand using multiple identical generators which can start, stop and operate efficiently at chosen load, independently of the others, making them suitable for baseload and peaking power generation. Matching supply and demand, called load balancing, is essential for a stable and reliable supply of electricity. Short-term deviations in the balance lead to frequency variations and a prolonged mismatch results in blackouts. Operators of power transmission systems a charged with the balancing task, matching the power output of all the generators to the load of their electrical grid. The load balancing task has become much more challenging as increasingly intermittent and variable generators such as wind turbines and solar cells are added to the grid, forcing other producers to adapt their output much more frequently than has been required in the past. The first two dynamic grid stability power plants utilizing the concept have been ordered by Elering and will be built by Wärtsilä in Kiisa, Estonia (Kiisa Power Plant). Their purpose is to "provide dynamic generation capacity to meet sudden and unexpected drops in the electricity supply". They are scheduled to be ready during 2013 and 2014, and their total output will be 250 MW. Power system automation enables rapid diagnosis of and precise solutions to specific grid disruptions or outages. These technologies rely on and contribute to each of the other four key areas. Three technology categories for advanced control methods are distributed intelligent agents (control systems), analytical tools (software algorithms and high-speed computers), and operational applications (SCADA, substation automation, demand response, etc.). Using artificial intelligence programming techniques, the Fujian power grid in China created a wide area protection system that is rapidly able to accurately calculate a control strategy and execute it. The Voltage Stability Monitoring & Control (VSMC) software uses a sensitivity-based successive linear programming method to reliably determine the optimal control solution. Research Major programs IntelliGrid – Created by the Electric Power Research Institute (EPRI), IntelliGrid architecture provides methodology, tools, and recommendations for standards and technologies for utility use in planning, specifying, and procuring IT-based systems, such as advanced metering, distribution automation, and demand response. The architecture also provides a living laboratory for assessing devices, systems, and technology. Several utilities have applied IntelliGrid architecture including Southern California Edison, Long Island Power Authority, Salt River Project, and TXU Electric Delivery. The IntelliGrid Consortium is a public/private partnership that integrates and optimizes global research efforts, funds technology R&D, works to integrate technologies, and disseminates technical information. Grid 2030 – Grid 2030 is a joint vision statement for the U.S. electrical system developed by the electric utility industry, equipment manufacturers, information technology providers, federal and state government agencies, interest groups, universities, and national laboratories. It covers generation, transmission, distribution, storage, and end-use. The National Electric Delivery Technologies Roadmap is the implementation document for the Grid 2030 vision. The Roadmap outlines the key issues and challenges for modernizing the grid and suggests paths that government and industry can take to build America's future electric delivery system. Modern Grid Initiative (MGI) is a collaborative effort between the U.S. Department of Energy (DOE), the National Energy Technology Laboratory (NETL), utilities, consumers, researchers, and other grid stakeholders to modernize and integrate the U.S. electrical grid. DOE's Office of Electricity Delivery and Energy Reliability (OE) sponsors the initiative, which builds upon Grid 2030 and the National Electricity Delivery Technologies Roadmap and is aligned with other programs such as GridWise and GridWorks. GridWise – A DOE OE program focused on developing information technology to modernize the U.S. electrical grid. Working with the GridWise Alliance, the program invests in communications architecture and standards; simulation and analysis tools; smart technologies; test beds and demonstration projects; and new regulatory, institutional, and market frameworks. The GridWise Alliance is a consortium of public and private electricity sector stakeholders, providing a forum for idea exchanges, cooperative efforts, and meetings with policy makers at federal and state levels. GridWise Architecture Council (GWAC) was formed by the U.S. Department of Energy to promote and enable interoperability among the many entities that interact with the nation's electric power system. The GWAC members are a balanced and respected team representing the many constituencies of the electricity supply chain and users. The GWAC provides industry guidance and tools to articulate the goal of interoperability across the electric system, identify the concepts and architectures needed to make interoperability possible, and develop actionable steps to facilitate the inter operation of the systems, devices, and institutions that encompass the nation's electric system. The GridWise Architecture Council Interoperability Context Setting Framework, V 1.1 defines necessary guidelines and principles. GridWorks – A DOE OE program focused on improving the reliability of the electric system through modernizing key grid components such as cables and conductors, substations and protective systems, and power electronics. The program's focus includes coordinating efforts on high temperature superconducting systems, transmission reliability technologies, electric distribution technologies, energy storage devices, and GridWise systems. Pacific Northwest Smart Grid Demonstration Project. - This project is a demonstration across five Pacific Northwest states-Idaho, Montana, Oregon, Washington, and Wyoming. It involves about 60,000 metered customers, and contains many key functions of the future smart grid. Solar Cities - In Australia, the Solar Cities programme included close collaboration with energy companies to trial smart meters, peak and off-peak pricing, remote switching and related efforts. It also provided some limited funding for grid upgrades. Smart Grid Energy Research Center (SMERC) - Located at University of California, Los Angeles dedicated its efforts to large-scale testing of its smart EV charging network technology. It created another platform for bidirectional flow of information between a utility and consumer end-devices. SMERC also developed a demand response (DR) test bed that comprises a Control Center, Demand Response Automation Server (DRAS), Home-Area-Network (HAN), Battery Energy Storage System (BESS), and photovoltaic (PV) panels. These technologies are installed within the Los Angeles Department of Water and Power and Southern California Edison territory as a network of EV chargers, battery energy storage systems, solar panels, DC fast charger, and Vehicle-to-Grid (V2G) units. These platforms, communications and control networks enables UCLA-led projects within the area to be tested in partnership with two local utilities, SCE and LADWP. Smart Quart - In Germany, the Smart Quart project develops three smart districts to develop, test and showcase technology to operate smart grids. The project is a collaboration of E.ON, Viessmann, gridX and hydrogenious together with the RWTH Aachen University. It is planned that by the end of 2024 all three districts are supplied with locally generated energy and are largely independent of fossil energy sources. Smart5Grid – In Portugal, aims to ensure that operators in the energy sector take advantage of the benefits associated with the use of 5G networks. With reliability and security, a solution is proposed to precisely address the specific requirements imposed by Smart Grids, such as high data transfer rates and real-time monitoring. Smart grid modelling Many different concepts have been used to model intelligent power grids. They are generally studied within the framework of complex systems. In a recent brainstorming session, the power grid was considered within the context of optimal control, ecology, human cognition, glassy dynamics, information theory, microphysics of clouds, and many others. Here is a selection of the types of analyses that have appeared in recent years. Protection systems that verify and supervise themselves Pelqim Spahiu and Ian R. Evans in their study introduced the concept of a substation based smart protection and hybrid Inspection Unit. Kuramoto oscillators The Kuramoto model is a well-studied system. The power grid has been described in this context as well. The goal is to keep the system in balance, or to maintain phase synchronization (also known as phase locking). Non-uniform oscillators also help to model different technologies, different types of power generators, patterns of consumption, and so on. The model has also been used to describe the synchronization patterns in the blinking of fireflies. Smart Grid Communication Network Network Simulators are used to simulate/emulate network communication effects. This typically involves setting up a lab with the smart grid devices, applications etc. with the virtual network being provided by the network simulator. Neural networks Neural networks have been considered for power grid management as well. Electric power systems can be classified in multiple different ways: non-linear, dynamic, discrete, or random. Artificial Neural Networks (ANNs) attempt to solve the most difficult of these problems, the non-linear problems. Demand Forecasting One application of ANNs is in demand forecasting. In order for grids to operate economically and reliably, demand forecasting is essential, because it is used to predict the amount of power that will be consumed by the load. This is dependent on weather conditions, type of day, random events, incidents, etc. For non-linear loads though, the load profile isn't smooth and as predictable, resulting in higher uncertainty and less accuracy using the traditional Artificial Intelligence models. Some factors that ANNs consider when developing these sort of models: classification of load profiles of different customer classes based on the consumption of electricity, increased responsiveness of demand to predict real time electricity prices as compared to conventional grids, the need to input past demand as different components, such as peak load, base load, valley load, average load, etc. instead of joining them into a single input, and lastly, the dependence of the type on specific input variables. An example of the last case would be given the type of day, whether its weekday or weekend, that wouldn't have much of an effect on Hospital grids, but it'd be a big factor in resident housing grids' load profile. Markov processes As wind power continues to gain popularity, it becomes a necessary ingredient in realistic power grid studies. Off-line storage, wind variability, supply, demand, pricing, and other factors can be modelled as a mathematical game. Here the goal is to develop a winning strategy. Markov processes have been used to model and study this type of system. Economics Market outlook In 2009, the US smart grid industry was valued at about $21.4 billion – by 2014, it will exceed at least $42.8 billion. Given the success of the smart grids in the U.S., the world market is expected to grow at a faster rate, surging from $69.3 billion in 2009 to $171.4 billion by 2014. With the segments set to benefit the most will be smart metering hardware sellers and makers of software used to transmit and organize the massive amount of data collected by meters. A 2011 study from the Electric Power Research Institute concludes that investment in a U.S. smart grid will cost up to $476 billion over 20 years but will provide up to $2 trillion in customer benefits over that time. In 2015, the World Economic Forum reported a transformational investment of more than $7.6 trillion by members of the OECD is needed over the next 25 years (or $300 billion per year) to modernize, expand, and decentralize the electricity infrastructure with technical innovation as key to the transformation. A 2019 study from International Energy Agency estimates that the current (depreciated) value of the US electric grid is more than USD 1 trillion. The total cost of replacing it with a smart grid is estimated to be more than USD 4 trillion. If smart grids are deployed fully across the US, the country expects to save USD 130 billion annually. General economics developments As customers can choose their electricity suppliers, depending on their different tariff methods, the focus of transportation costs will be increased. Reduction of maintenance and replacements costs will stimulate more advanced control. A smart grid precisely limits electrical power down to the residential level, network small-scale distributed energy generation and storage devices, communicate information on operating status and needs, collect information on prices and grid conditions, and move the grid beyond central control to a collaborative network. US and UK savings estimates and concerns A 2003 United States Department of Energy study calculated that internal modernization of US grids with smart grid capabilities would save between 46 and 117 billion dollars over the next 20 years if implemented within a few years of the study. As well as these industrial modernization benefits, smart grid features could expand energy efficiency beyond the grid into the home by coordinating low priority home devices such as water heaters so that their use of power takes advantage of the most desirable energy sources. Smart grids can also coordinate the production of power from large numbers of small power producers such as owners of rooftop solar panels — an arrangement that would otherwise prove problematic for power systems operators at local utilities. One important question is whether consumers will act in response to market signals. The U.S. Department of Energy (DOE) as part of the American Recovery and Reinvestment Act Smart Grid Investment Grant and Demonstrations Program funded special consumer behavior studies to examine the acceptance, retention, and response of consumers subscribed to time-based utility rate programs that involve advanced metering infrastructure and customer systems such as in-home displays and programmable communicating thermostats. Another concern is that the cost of telecommunications to fully support smart grids may be prohibitive. A less expensive communication mechanism is proposed using a form of "dynamic demand management" where devices shave peaks by shifting their loads in reaction to grid frequency. Grid frequency could be used to communicate load information without the need of an additional telecommunication network, but it would not support economic bargaining or quantification of contributions. Although there are specific and proven smart grid technologies in use, smart grid is an aggregate term for a set of related technologies on which a specification is generally agreed, rather than a name for a specific technology. Some of the benefits of such a modernized electricity network include the ability to reduce power consumption at the consumer side during peak hours, called demand side management; enabling grid connection of distributed generation power (with photovoltaic arrays, small wind turbines, micro hydro, or even combined heat power generators in buildings); incorporating grid energy storage for distributed generation load balancing; and eliminating or containing failures such as widespread power grid cascading failures. The increased efficiency and reliability of the smart grid is expected to save consumers money and help reduce emissions. Oppositions and concerns Most opposition and concerns have centered on smart meters and the items (such as remote control, remote disconnect, and variable rate pricing) enabled by them. Where opposition to smart meters is encountered, they are often marketed as "smart grid" which connects smart grid to smart meters in the eyes of opponents. Specific points of opposition or concern include: consumer concerns over privacy, e.g. use of usage data by law enforcement social concerns over "fair" availability of electricity concern that complex rate systems (e.g. variable rates) remove clarity and accountability, allowing the supplier to take advantage of the customer concern over remotely controllable "kill switch" incorporated into most smart meters social concerns over Enron style abuses of information leverage concerns over giving the government mechanisms to control the use of all power using activities concerns over RF emissions from smart meters Security While modernization of electrical grids into smart grids allows for optimization of everyday processes, a smart grid, being online, can be vulnerable to cyberattacks. Transformers which increase the voltage of electricity created at power plants for long-distance travel, transmission lines themselves, and distribution lines which deliver the electricity to its consumers are particularly susceptible. These systems rely on sensors which gather information from the field and then deliver it to control centers, where algorithms automate analysis and decision-making processes. These decisions are sent back to the field, where existing equipment execute them. Hackers have the potential to disrupt these automated control systems, severing the channels which allow generated electricity to be utilized. This is called a denial of service or DoS attack. They can also launch integrity attacks which corrupt information being transmitted along the system as well as desynchronization attacks which affect when such information is delivered to the appropriate location. Additionally, intruders can gain access via renewable energy generation systems and smart meters connected to the grid, taking advantage of more specialized weaknesses or ones whose security has not been prioritized. Because a smart grid has a large number of access points, like smart meters, defending all of its weak points can prove difficult. There is also concern on the security of the infrastructure, primarily that involving communications technology. Concerns chiefly center around the communications technology at the heart of the smart grid. Designed to allow real-time contact between utilities and meters in customers' homes and businesses, there is a risk that these capabilities could be exploited for criminal or even terrorist actions. One of the key capabilities of this connectivity is the ability to remotely switch off power supplies, enabling utilities to quickly and easily cease or modify supplies to customers who default on payment. This is undoubtedly a massive boon for energy providers, but also raises some significant security issues. Cybercriminals have infiltrated the U.S. electric grid before on numerous occasions. Aside from computer infiltration, there are also concerns that computer malware like Stuxnet, which targeted SCADA systems which are widely used in industry, could be used to attack a smart grid network. Electricity theft is a concern in the U.S. where the smart meters being deployed use RF technology to communicate with the electricity transmission network. People with knowledge of electronics can devise interference devices to cause the smart meter to report lower than actual usage. Similarly, the same technology can be employed to make it appear that the energy the consumer is using is being used by another customer, increasing their bill. The damage from a well-executed, sizable cyberattack could be extensive and long-lasting. One incapacitated substation could take from nine days to over a year to repair, depending on the nature of the attack. It can also cause an hours-long outage in a small radius. It could have an immediate effect on transportation infrastructure, as traffic lights and other routing mechanisms as well as ventilation equipment for underground roadways is reliant on electricity. Additionally, infrastructure which relies on the electric grid, including wastewater treatment facilities, the information technology sector, and communications systems could be impacted. The December 2015 Ukraine power grid cyberattack, the first recorded of its kind, disrupted services to nearly a quarter of a million people by bringing substations offline. The Council on Foreign Relations has noted that states are most likely to be the perpetrators of such an attack as they have access to the resources to carry one out despite the high level of difficulty of doing so. Cyber intrusions can be used as portions of a larger offensive, military or otherwise. Some security experts warn that this type of event is easily scalable to grids elsewhere. Insurance company Lloyd's of London has already modeled the outcome of a cyberattack on the Eastern Interconnection, which has the potential to impact 15 states, put 93 million people in the dark, and cost the country's economy anywhere from $243 billion to $1 trillion in various damages. According to the U.S. House of Representatives Subcommittee on Economic Development, Public Buildings, and Emergency Management, the electric grid has already seen a sizable number of cyber intrusions, with two in every five aiming to incapacitate it. As such, the U.S. Department of Energy has prioritized research and development to decrease the electric grid's vulnerability to cyberattacks, citing them as an "imminent danger" in its 2017 Quadrennial Energy Review. The Department of Energy has also identified both attack resistance and self-healing as major keys to ensuring that today's smart grid is future-proof. While there are regulations already in place, namely the Critical Infrastructure Protection Standards introduced by the North America Electric Reliability Council, a significant number of them are suggestions rather than mandates. Most electricity generation, transmission, and distribution facilities and equipment are owned by private stakeholders, further complicating the task of assessing adherence to such standards. Additionally, even if utilities want to fully comply, they may find that it is too expensive to do so. Some experts argue that the first step to increasing the cyber defenses of the smart electric grid is completing a comprehensive risk analysis of existing infrastructure, including research of software, hardware, and communication processes. Additionally, as intrusions themselves can provide valuable information, it could be useful to analyze system logs and other records of their nature and timing. Common weaknesses already identified using such methods by the Department of Homeland Security include poor code quality, improper authentication, and weak firewall rules. Once this step is completed, some suggest that it makes sense to then complete an analysis of the potential consequences of the aforementioned failures or shortcomings. This includes both immediate consequences as well as second- and third-order cascading effects on parallel systems. Finally, risk mitigation solutions, which may include simple remediation of infrastructure inadequacies or novel strategies, can be deployed to address the situation. Some such measures include recoding of control system algorithms to make them more able to resist and recover from cyberattacks or preventive techniques that allow more efficient detection of unusual or unauthorized changes to data. Strategies to account for human error which can compromise systems include educating those who work in the field to be wary of strange USB drives, which can introduce malware if inserted, even if just to check their contents. Other solutions include utilizing transmission substations, constrained SCADA networks, policy based data sharing, and attestation for constrained smart meters. Transmission substations utilize one-time signature authentication technologies and one-way hash chain constructs. These constraints have since been remedied with the creation of a fast-signing and verification technology and buffering-free data processing. A similar solution has been constructed for constrained SCADA networks. This involves applying a Hash-Based Message Authentication Code to byte streams, converting the random-error detection available on legacy systems to a mechanism that guarantees data authenticity. Policy-based data sharing utilizes GPS-clock-synchronized-fine-grain power grid measurements to provide increased grid stability and reliability. It does this through synchro-phasor requirements that are gathered by PMUs. Attestation for constrained smart meters faces a slightly different challenge, however. One of the biggest issues with attestation for constrained smart meters is that in order to prevent energy theft, and similar attacks, cyber security providers have to make sure that the devices' software is authentic. To combat this problem, an architecture for constrained smart networks has been created and implemented at a low level in the embedded system. The protection system of a smart grid provides grid reliability analysis, failure protection, and security and privacy protection services. While the additional communication infrastructure of a smart grid provides additional protective and security mechanisms, it also presents a risk of external attack and internal failures. In a report on cyber security of smart grid technology first produced in 2010, and later updated in 2014, the US National Institute of Standards and Technology pointed out that the ability to collect more data about energy use from customer smart meters also raises major privacy concerns, since the information stored at the meter, which is potentially vulnerable to data breaches, can be mined for personal details about customers. Other challenges to adoption Before a utility installs an advanced metering system, or any type of smart system, it must make a business case for the investment. Some components, like the power system stabilizers (PSS) installed on generators are very expensive, require complex integration in the grid's control system, are needed only during emergencies, and are only effective if other suppliers on the network have them. Without any incentive to install them, power suppliers don't. Most utilities find it difficult to justify installing a communications infrastructure for a single application (e.g. meter reading). Because of this, a utility must typically identify several applications that will use the same communications infrastructure – for example, reading a meter, monitoring power quality, remote connection and disconnection of customers, enabling demand response, etc. Ideally, the communications infrastructure will not only support near-term applications, but unanticipated applications that will arise in the future. Regulatory or legislative actions can also drive utilities to implement pieces of a smart grid puzzle. Each utility has a unique set of business, regulatory, and legislative drivers that guide its investments. This means that each utility will take a different path to creating their smart grid and that different utilities will create smart grids at different adoption rates. Some features of smart grids draw opposition from industries that currently are, or hope to provide similar services. An example is competition with cable and DSL Internet providers from broadband over powerline internet access. Providers of SCADA control systems for grids have intentionally designed proprietary hardware, protocols and software so that they cannot inter-operate with other systems in order to tie its customers to the vendor. The incorporation of digital communications and computer infrastructure with the grid's existing physical infrastructure poses challenges and inherent vulnerabilities. According to IEEE Security and Privacy Magazine, the smart grid will require that people develop and use large computer and communication infrastructure that supports a greater degree of situational awareness and that allows for more specific command and control operations. This process is necessary to support major systems such as demand-response wide-area measurement and control, storage and transportation of electricity, and the automation of electric distribution. Power Theft / Power Loss Various "smart grid" systems have dual functions. This includes Advanced Metering Infrastructure systems which, when used with various software can be used to detect power theft and by process of elimination, detect where equipment failures have taken place. These are in addition to their primary functions of eliminating the need for human meter reading and measuring the time-of-use of electricity. The worldwide power loss including theft is estimated at two-hundred billion dollars annually. Electricity theft also represents a major challenge when providing reliable electrical service in developing countries. Deployments and attempted deployments Enel The earliest, and one of the largest, example of a smart grid is the Italian system installed by Enel S.p.A. of Italy. Completed in 2005, the Telegestore project was highly unusual in the utility world because the company designed and manufactured their own meters, acted as their own system integrator, and developed their own system software. The Telegestore project is widely regarded as the first commercial scale use of smart grid technology to the home, and delivers annual savings of 500 million euro at a project cost of 2.1 billion euro. US Dept. of Energy - ARRA Smart Grid Project One of the largest deployment programs in the world to-date is the U.S. Dept. of Energy's Smart Grid Program funded by the American Recovery and Reinvestment Act of 2009. This program required matching funding from individual utilities. A total of over $9 billion in Public/Private funds were invested as part of this program. Technologies included Advanced Metering Infrastructure, including over 65 million Advanced "Smart" Meters, Customer Interface Systems, Distribution & Substation Automation, Volt/VAR Optimization Systems, over 1,000 Synchrophasors, Dynamic Line Rating, Cyber Security Projects, Advanced Distribution Management Systems, Energy Storage Systems, and Renewable Energy Integration Projects. This program consisted of Investment Grants (matching), Demonstration Projects, Consumer Acceptance Studies, and Workforce Education Programs. Reports from all individual utility programs as well as overall impact reports will be completed by the second quarter of 2015. In the U.S., the Energy Policy Act of 2005 and Title XIII of the Energy Independence and Security Act of 2007 are providing funding to encourage smart grid development. The objective is to enable utilities to better predict their needs, and in some cases involve consumers in a time-of-use tariff. Funds have also been allocated to develop more robust energy control technologies. Austin, Texas In the US, the city of Austin, Texas, has been working on building its smart grid since 2003, when its utility first replaced 1/3 of its manual meters with smart meters that communicate via a wireless mesh network. It currently manages 200,000 devices real-time (smart meters, smart thermostats, and sensors across its service area), and expects to be supporting 500,000 devices real-time in 2009 servicing 1 million consumers and 43,000 businesses. Boulder, Colorado Boulder, Colorado, completed the first phase of its smart grid project in August 2008. Both systems use the smart meter as a gateway to the home automation network (HAN) that controls smart sockets and devices. Some HAN designers favor decoupling control functions from the meter, out of concern of future mismatches with new standards and technologies available from the fast moving business segment of home electronic devices. Hydro One Hydro One, in Ontario, Canada is in the midst of a large-scale Smart Grid initiative, deploying a standards-compliant communications infrastructure from Trilliant. By the end of 2010, the system will serve 1.3 million customers in the province of Ontario. The initiative won the "Best AMR Initiative in North America" award from the Utility Planning Network. Île d'Yeu Île d'Yeu began a 2-year pilot program in Spring of 2020. Twenty-three houses in the Ker Pissot neighborhood and surrounding areas were interconnected with a microgrid that was automated as a smart grid with software from Engie. Sixty-four solar panels with a peak capacity of 23.7 kW were installed on five houses and a battery with a storage capacity of 15 kWh was installed on one house. Six houses store excess solar energy in their hot water heaters. A dynamic system apportions the energy provided by the solar panels and stored in the battery and hot water heaters to the system of 23 houses. The smart grid software dynamically updates energy supply and demand in 5 minute intervals, deciding whether to pull energy from the battery or from the panels and when to store it in the hot water heaters. This pilot program was the first such project in France. Mannheim The City of Mannheim in Germany is using realtime Broadband Powerline (BPL) communications in its Model City Mannheim "MoMa" project. Sydney Sydney also in Australia, in partnership with the Australian Government implemented the Smart Grid, Smart City program. Évora InovGrid is an innovative project in Évora, Portugal that aims to equip the electricity grid with information and devices to automate grid management, improve service quality, reduce operating costs, promote energy efficiency and environmental sustainability, and increase the penetration of renewable energies and electric vehicles. It will be possible to control and manage the state of the entire electricity distribution grid at any given instant, allowing suppliers and energy services companies to use this technological platform to offer consumers information and added-value energy products and services. This project to install an intelligent energy grid places Portugal and EDP at the cutting edge of technological innovation and service provision in Europe. E-Energy In the so-called E-Energy projects several German utilities are creating first nucleolus in six independent model regions. A technology competition identified this model regions to carry out research and development activities with the main objective to create an "Internet of Energy." Massachusetts One of the first attempted deployments of "smart grid" technologies in the United States was rejected in 2009 by electricity regulators in the Commonwealth of Massachusetts, a US state. According to an article in the Boston Globe, Northeast Utilities' Western Massachusetts Electric Co. subsidiary actually attempted to create a "smart grid" program using public subsidies that would switch low income customers from post-pay to pre-pay billing (using "smart cards") in addition to special hiked "premium" rates for electricity used above a predetermined amount. This plan was rejected by regulators as it "eroded important protections for low-income customers against shutoffs". According to the Boston Globe, the plan "unfairly targeted low-income customers and circumvented Massachusetts laws meant to help struggling consumers keep the lights on". A spokesman for an environmental group supportive of smart grid plans and Western Massachusetts' Electric's aforementioned "smart grid" plan, in particular, stated "If used properly, smart grid technology has a lot of potential for reducing peak demand, which would allow us to shut down some of the oldest, dirtiest power plants... It's a tool." eEnergy Vermont consortium The eEnergy Vermont consortium is a US statewide initiative in Vermont, funded in part through the American Recovery and Reinvestment Act of 2009, in which all of the electric utilities in the state have rapidly adopted a variety of Smart Grid technologies, including about 90% Advanced Metering Infrastructure deployment, and are presently evaluating a variety of dynamic rate structures. Netherlands In the Netherlands a large-scale project (>5000 connections, >20 partners) was initiated to demonstrate integrated smart grids technologies, services and business cases. Chattanooga EPB in Chattanooga, TN is a municipally-owned electric utility that started construction of a smart grid in 2008, receiving a $111,567,606 grant from the US DOE in 2009 to expedite construction and implementation (for a total budget of $232,219,350). Deployment of power-line interrupters (1170 units) was completed in April 2012, and deployment of smart meters (172,079 units) was completed in 2013. The smart grid's backbone fiber-optic system was also used to provide the first gigabit-speed internet connection to residential customers in the US through the Fiber to the Home initiative, and now speeds of up to 10 gigabits per second are available to residents. The smart grid is estimated to have reduced power outages by an average of 60%, saving the city about 60 million dollars annually. It has also reduced the need for "truck rolls" to scout and troubleshoot faults, resulting in an estimated reduction of 630,000 truck driving miles, and 4.7 million pounds of carbon emissions. In January 2016, EPB became the first major power distribution system to earn Performance Excellence in Electricity Renewal (PEER) certification. OpenADR Implementations Certain deployments utilize the OpenADR standard for load shedding and demand reduction during higher demand periods. China The smart grid market in China is estimated to be $22.3 billion with a projected growth to $61.4 billion by 2015. Honeywell is developing a demand response pilot and feasibility study for China with the State Grid Corp. of China using the OpenADR demand response standard. The State Grid Corp., the Chinese Academy of Science, and General Electric intend to work together to develop standards for China's smart grid rollout. United States In 2009, the US Department of Energy awarded an $11 million grant to Southern California Edison and Honeywell for a demand response program that automatically turns down energy use during peak hours for participating industrial customers. The Department of Energy awarded an $11.4 million grant to Honeywell to implement the program using the OpenADR standard. Hawaiian Electric Co. (HECO) is implementing a two-year pilot project to test the ability of an ADR program to respond to the intermittence of wind power. Hawaii has a goal to obtain 70 percent of its power from renewable sources by 2030. HECO will give customers incentives for reducing power consumption within 10 minutes of a notice. Guidelines, standards and user groups Part of the IEEE Smart Grid Initiative, IEEE 2030.2 represents an extension of the work aimed at utility storage systems for transmission and distribution networks. The IEEE P2030 group expects to deliver early 2011 an overarching set of guidelines on smart grid interfaces. The new guidelines will cover areas including batteries and supercapacitors as well as flywheels. The group has also spun out a 2030.1 effort drafting guidelines for integrating electric vehicles into the smart grid. IEC TC 57 has created a family of international standards that can be used as part of the smart grid. These standards include IEC 61850 which is an architecture for substation automation, and IEC 61970/61968 – the Common Information Model (CIM). The CIM provides for common semantics to be used for turning data into information. OpenADR is an open-source smart grid communications standard used for demand response applications. It is typically used to send information and signals to cause electrical power-using devices to be turned off during periods of higher demand. MultiSpeak has created a specification that supports distribution functionality of the smart grid. MultiSpeak has a robust set of integration definitions that supports nearly all of the software interfaces necessary for a distribution utility or for the distribution portion of a vertically integrated utility. MultiSpeak integration is defined using extensible markup language (XML) and web services. The IEEE has created a standard to support synchrophasors – C37.118. The UCA International User Group discusses and supports real world experience of the standards used in smart grids. A utility task group within LonMark International deals with smart grid related issues. There is a growing trend towards the use of TCP/IP technology as a common communication platform for smart meter applications, so that utilities can deploy multiple communication systems, while using IP technology as a common management platform. IEEE P2030 is an IEEE project developing a "Draft Guide for Smart Grid Interoperability of Energy Technology and Information Technology Operation with the Electric Power System (EPS), and End-Use Applications and Loads". NIST has included ITU-T G.hn as one of the "Standards Identified for Implementation" for the Smart Grid "for which it believed there was strong stakeholder consensus". G.hn is standard for high-speed communications over power lines, phone lines and coaxial cables. OASIS EnergyInterop' – An OASIS technical committee developing XML standards for energy interoperation. Its starting point is the California OpenADR standard. Under the Energy Independence and Security Act of 2007 (EISA), NIST is charged with overseeing the identification and selection of hundreds of standards that will be required to implement the Smart Grid in the U.S. These standards will be referred by NIST to the Federal Energy Regulatory Commission (FERC). This work has begun, and the first standards have already been selected for inclusion in NIST's Smart Grid catalog. However, some commentators have suggested that the benefits that could be realized from Smart Grid standardization could be threatened by a growing number of patents that cover Smart Grid architecture and technologies. If patents that cover standardized Smart Grid elements are not revealed until technology is broadly distributed throughout the network ("locked-in"), significant disruption could occur when patent holders seek to collect unanticipated rents from large segments of the market. GridWise Alliance rankings In November 2017 the non-profit GridWise Alliance along with Clean Edge Inc., a clean energy group, released rankings for all 50 states in their efforts to modernize the electric grid. California was ranked number one. The other top states were Illinois, Texas, Maryland, Oregon, Arizona, the District of Columbia, New York, Nevada and Delaware. "The 30-plus page report from the GridWise Alliance, which represents stakeholders that design, build and operate the electric grid, takes a deep dive into grid modernization efforts across the country and ranks them by state."	Technology	Electricity transmission and distribution	null
9214162	https://en.wikipedia.org/wiki/Turritopsis%20dohrnii	Turritopsis dohrnii	Turritopsis dohrnii, also known as the immortal jellyfish, is a species of small, biologically immortal jellyfish found worldwide in temperate to tropic waters. It is one of the few known cases of animals capable of reverting completely to a sexually immature, colonial stage after having reached sexual maturity as a solitary individual. Like most other hydrozoans, T. dohrnii begin their lives as tiny, free-swimming larvae known as planulae. As a planula settles down, it gives rise to a colony of polyps that are attached to the sea floor. All the polyps and jellyfish arising from a single planula are genetically identical clones. The polyps form into an extensively branched form, which is not commonly seen in most jellyfish. Jellyfish, also known as medusae, then bud off these polyps and continue their life in a free-swimming form, eventually becoming sexually mature. When sexually mature, they have been known to prey on other jellyfish species at a rapid pace. If the T. dohrnii jellyfish is exposed to environmental stress, physical assault, or is sick or old, it can revert to the polyp stage, forming a new polyp colony. It does this through the cell development process of transdifferentiation, which alters the differentiated state of the cells and transforms them into new types of cells. Theoretically, this process can go on indefinitely, effectively rendering the jellyfish biologically immortal, although in practice individuals can still die. In nature, most Turritopsis dohrnii are likely to succumb to predation or disease in the medusa stage without reverting to the polyp form. The capability of biological immortality with no maximum lifespan makes T. dohrnii an important target of basic biological aging and pharmaceutical research. Taxonomy The species was formerly considered conspecific with T. nutricula before being reclassified as a separate species. It was named in 1883 in honour of Anton Dohrn, the founder of the Stazione Zoologica Anton Dohrn in Naples, Italy. Until a 2006 study, it was thought that Turritopsis rubra and Turritopsis nutricula were the same species as Turritopsis dohrnii. It is not known whether or not T. rubra medusae can also transform back into polyps, however further research is still to be done. Description The medusa of Turritopsis dohrnii is bell-shaped, with a maximum diameter of about and is about as tall as it is wide. The mesoglea in the walls of the bell is uniformly thin, except for some thickening at the apex. The relatively large stomach is bright red and has a cruciform shape in cross section. Young specimens 1 mm in diameter have only eight tentacles evenly spaced out along the edge, whereas adult specimens have 80–90 tentacles. The medusa (jellyfish) is free-living in the plankton. Dense nerve net cells are also present in the epidermis in the cap. They form a large ring-like structure above the radial canal commonly presented in cnidarians. Turritopsis dohrnii also has a bottom-living polyp form, or hydroid, which consists of stolons that run along the substrate and upright branches with feeding polyps that can produce medusa buds. These polyps develop over a few days into tiny 1 mm medusae, which are liberated and swim free from the parent hydroid colony. Distribution and invasion Turritopsis is believed to have originated in the Pacific, but has spread all over the world through trans-Arctic migrations, and has speciated into several populations that are easy to distinguish morphologically, but whose species distinctions have recently been verified by a study and comparison of mitochondrial ribosomal gene sequences. Turritopsis are found in temperate to tropical regions in all of the world's oceans. Turritopsis is believed to be spreading across the world through ballast water discharge. Unlike other species invasions which caused serious economic and ecological consequences, T. dohrnii's invasion around the world was unnoticed due to their tiny size and innocuity. "We are looking at a worldwide silent invasion", said Smithsonian Tropical Marine Institute scientist Maria Miglietta. Life cycle The eggs develop in gonads of female medusae, which are located in the walls of the manubrium (stomach). Mature eggs are presumably spawned and fertilized in the sea by sperm produced and released by male medusae, as is the case for most hydromedusae. However, the related species Turritopsis rubra seems to retain fertilized eggs until the planula stage. Fertilized eggs develop into planula larvae, which settle onto the sea floor (or even the rich marine communities that live on floating docks), and develop into polyp colonies (hydroids). The hydroids bud new jellyfishes, which are released at about one millimetre in size and then grow and feed in the plankton, becoming sexually mature after a few weeks (the exact duration depends on the ocean temperature; at it is 25 to 30 days and at it is 18 to 22 days). Medusae of T. dohrnii are able to survive between 14 °C and 25 °C. Biological immortality Most jellyfish species have a relatively fixed lifespan, which varies by species from hours to many months (long-lived mature jellyfish spawn every day or night; the time is also fairly fixed and species-specific). The medusa of Turritopsis dohrnii is the only form known to have the ability to return to a polyp state, by a specific transformation process that requires the presence of certain cell types (tissue from both the jellyfish bell surface and the circulatory canal system). Experiments have revealed that all stages of the medusae, from newly released to fully mature individuals, can transform back into polyps under the conditions of starvation, sudden temperature change, reduction of salinity, and artificial damage of the bell with forceps or scissors. The transforming medusa is characterized first by deterioration of the bell, mesoglea, and tentacles. All immature medusa (with 12 tentacles at most) then turned into a cyst-like stage and then transformed into stolons and polyps. However, about 20%-40% of mature medusa went into the stolons and polyps stage without passing the cyst-like stage. Polyps were formed after 2 days since stolons had developed and fed on food. Polyps further multiply by growing additional stolons, branches, and then polyps to form colonial hydroids. In the experiment, they would eventually transform into stolons and polyps and begin their lives once again, even without environmental changes or injury. This ability to reverse the biotic cycle (in response to adverse conditions) is unique in the animal kingdom. It allows the jellyfish to bypass death, rendering Turritopsis dohrnii potentially biologically immortal. The process has not been observed in their natural habitat, in part because the process is quite rapid and because field observations at the right moment are unlikely. Regardless, most individual medusae are likely to fall victim to the general hazards of life as mesoplankton, including being eaten by predators or succumbing to disease. The species possesses unique mechanisms related to telomere maintenance, which play a significant role in its regenerative abilities. T. dohrnii maintains telomere length through specific cellular processes during its life cycle reversal, effectively resetting cellular aging. The species' cell development method of transdifferentiation has inspired scientists to find a way to make stem cells using this process for renewing damaged or dead tissue in humans. Ecology Diet Turritopsis dohrnii are a carnivorous species that commonly feed on zooplankton. Their diet mainly consists of plankton, fish eggs and small mollusks. T. dohrnii ingests food and excretes waste through the mouth. T. dohrnii hunts by using its tentacles as it drifts through the water. Its tentacles, which contain stinging cells called nematocysts, spread and sting its prey. The tentacles can then flex to direct its prey to the mouth. T. dohrnii, like other jellyfish, may use its bell to catch its prey. T. dohrnii's bell will expand, sucking in water, as it propels itself to swim. This expansion of the bell brings potential prey in closer reach of the tentacles. Predation Turritopsis dohrnii, like other jellyfish, are preyed on most commonly by other jellyfish. Other predators of T. dohrnii include sea anemones, tuna, sharks, swordfish, sea turtles, and penguins. Many species prey on T. dohrnii and other jellyfish due to their simple composition. They are only approximately 5% non-aqueous matter, and the remaining part is composed of water. They are composed of three layers. An outer layer (the epidermis), a middle layer (mesoglea; a thick, jelly-like substance), and an inner layer (gastrodermis). Habitat Turritopsis dohrnii was first discovered in the Mediterranean Sea, but has since been found worldwide. T. dohrnii is generally found living in temperate to tropical waters. They can be found in marinas or docks, on vessel hulls, and on the ocean floor. They typically live in a salinity range of polyhaline (18–30 PSU) and euhaline (30-40 PSU). Genomic analysis Genomic analyses such as sequence analysis on mRNA or mitochondria DNA have been employed to investigate its lifecycle. mRNA analysis of each life stage showed that a stage-specific gene in the medusae stage is expressed tenfold more than in other stages. This gene is relative to a Wnt signal that can induce a regeneration process upon injury. Analysis of nucleotide sequence homologs and protein homologs identified Nemopsis bachei as the species' closest relative. None of the closely related species display biological immortality. In 2022, a study reported the key molecular mechanisms of rejuvenation they found in a comparison of the newly presented genomes of this biologically immortal jellyfish and a similar but non-rejuvenating jellyfish, involving e.g. DNA replication and repair, and stem cell renewal. Culturing Keeping T. dohrnii in captivity is quite difficult. Currently, only one scientist, Shin Kubota from Kyoto University, has managed to sustain a group of these jellyfish for a prolonged period of time. The plankton must be inspected daily to ensure that they have properly digested the Artemia cysts they are being fed. Kubota reported that during a two-year period, his colony rebirthed itself 11 times. Kubota regularly appears on Japanese television to talk about his immortal jellyfish and has recorded several songs about them, often singing them at the end of his conference presentations.	Biology and health sciences	Cnidarians	Animals
9217017	https://en.wikipedia.org/wiki/Transition%20state%20theory	Transition state theory	In chemistry, transition state theory (TST) explains the reaction rates of elementary chemical reactions. The theory assumes a special type of chemical equilibrium (quasi-equilibrium) between reactants and activated transition state complexes. TST is used primarily to understand qualitatively how chemical reactions take place. TST has been less successful in its original goal of calculating absolute reaction rate constants because the calculation of absolute reaction rates requires precise knowledge of potential energy surfaces, but it has been successful in calculating the standard enthalpy of activation (ΔH‡, also written Δ‡Hɵ), the standard entropy of activation (ΔS‡ or Δ‡Sɵ), and the standard Gibbs energy of activation (ΔG‡ or Δ‡Gɵ) for a particular reaction if its rate constant has been experimentally determined (the ‡ notation refers to the value of interest at the transition state; ΔH‡ is the difference between the enthalpy of the transition state and that of the reactants). This theory was developed simultaneously in 1935 by Henry Eyring, then at Princeton University, and by Meredith Gwynne Evans and Michael Polanyi of the University of Manchester. TST is also referred to as "activated-complex theory", "absolute-rate theory", and "theory of absolute reaction rates". Before the development of TST, the Arrhenius rate law was widely used to determine energies for the reaction barrier. The Arrhenius equation derives from empirical observations and ignores any mechanistic considerations, such as whether one or more reactive intermediates are involved in the conversion of a reactant to a product. Therefore, further development was necessary to understand the two parameters associated with this law, the pre-exponential factor (A) and the activation energy (Ea). TST, which led to the Eyring equation, successfully addresses these two issues; however, 46 years elapsed between the publication of the Arrhenius rate law, in 1889, and the Eyring equation derived from TST, in 1935. During that period, many scientists and researchers contributed significantly to the development of the theory. Theory The basic ideas behind transition state theory are as follows: Rates of reaction can be studied by examining activated complexes near the saddle point of a potential energy surface. The details of how these complexes are formed are not important. The saddle point itself is called the transition state. The activated complexes are in a special equilibrium (quasi-equilibrium) with the reactant molecules. The activated complexes can convert into products, and kinetic theory can be used to calculate the rate of this conversion. Development In the development of TST, three approaches were taken as summarized below. Thermodynamic treatment In 1884, Jacobus van 't Hoff proposed the Van 't Hoff equation describing the temperature dependence of the equilibrium constant for a reversible reaction: {A} <=> {B} where ΔU is the change in internal energy, K is the equilibrium constant of the reaction, R is the universal gas constant, and T is thermodynamic temperature. Based on experimental work, in 1889, Svante Arrhenius proposed a similar expression for the rate constant of a reaction, given as follows: Integration of this expression leads to the Arrhenius equation where k is the rate constant. A was referred to as the frequency factor (now called the pre-exponential coefficient), and Ea is regarded as the activation energy. By the early 20th century many had accepted the Arrhenius equation, but the physical interpretation of A and Ea remained vague. This led many researchers in chemical kinetics to offer different theories of how chemical reactions occurred in an attempt to relate A and Ea to the molecular dynamics directly responsible for chemical reactions. In 1910, French chemist René Marcelin introduced the concept of standard Gibbs energy of activation. His relation can be written as At about the same time as Marcelin was working on his formulation, Dutch chemists Philip Abraham Kohnstamm, Frans Eppo Cornelis Scheffer, and Wiedold Frans Brandsma introduced standard entropy of activation and the standard enthalpy of activation. They proposed the following rate constant equation However, the nature of the constant was still unclear. Kinetic-theory treatment In early 1900, Max Trautz and William Lewis studied the rate of the reaction using collision theory, based on the kinetic theory of gases. Collision theory treats reacting molecules as hard spheres colliding with one another; this theory neglects entropy changes, since it assumes that the collision between molecules are completely elastic. Lewis applied his treatment to the following reaction and obtained good agreement with experimental result. 2 HI → H2 + I2 However, later when the same treatment was applied to other reactions, there were large discrepancies between theoretical and experimental results. Statistical-mechanical treatment Statistical mechanics played a significant role in the development of TST. However, the application of statistical mechanics to TST was developed very slowly given the fact that in mid-19th century, James Clerk Maxwell, Ludwig Boltzmann, and Leopold Pfaundler published several papers discussing reaction equilibrium and rates in terms of molecular motions and the statistical distribution of molecular speeds. It was not until 1912 when the French chemist A. Berthoud used the Maxwell–Boltzmann distribution law to obtain an expression for the rate constant. where a and b are constants related to energy terms. Two years later, René Marcelin made an essential contribution by treating the progress of a chemical reaction as a motion of a point in phase space. He then applied Gibbs' statistical-mechanical procedures and obtained an expression similar to the one he had obtained earlier from thermodynamic consideration. In 1915, another important contribution came from British physicist James Rice. Based on his statistical analysis, he concluded that the rate constant is proportional to the "critical increment". His ideas were further developed by Richard Chace Tolman. In 1919, Austrian physicist Karl Ferdinand Herzfeld applied statistical mechanics to the equilibrium constant and kinetic theory to the rate constant of the reverse reaction, k−1, for the reversible dissociation of a diatomic molecule. AB <=>[k_1][k_{-1}] {A} + {B} He obtained the following equation for the rate constant of the forward reaction where is the dissociation energy at absolute zero, kB is the Boltzmann constant, h is the Planck constant, T is thermodynamic temperature, is vibrational frequency of the bond. This expression is very important since it is the first time that the factor kBT/h, which is a critical component of TST, has appeared in a rate equation. In 1920, the American chemist Richard Chace Tolman further developed Rice's idea of the critical increment. He concluded that critical increment (now referred to as activation energy) of a reaction is equal to the average energy of all molecules undergoing reaction minus the average energy of all reactant molecules. Potential energy surfaces The concept of potential energy surface was very important in the development of TST. The foundation of this concept was laid by René Marcelin in 1913. He theorized that the progress of a chemical reaction could be described as a point in a potential energy surface with coordinates in atomic momenta and distances. In 1931, Henry Eyring and Michael Polanyi constructed a potential energy surface for the reaction below. This surface is a three-dimensional diagram based on quantum-mechanical principles as well as experimental data on vibrational frequencies and energies of dissociation. H + H2 → H2 + H A year after the Eyring and Polanyi construction, Hans Pelzer and Eugene Wigner made an important contribution by following the progress of a reaction on a potential energy surface. The importance of this work was that it was the first time that the concept of col or saddle point in the potential energy surface was discussed. They concluded that the rate of a reaction is determined by the motion of the system through that col. Kramers theory of reaction rates By modeling reactions as Langevin motion along a one dimensional reaction coordinate, Hendrik Kramers was able to derive a relationship between the shape of the potential energy surface along the reaction coordinate and the transition rates of the system. The formulation relies on approximating the potential energy landscape as a series of harmonic wells. In a two state system, there will be three wells; a well for state A, an upside-down well representing the potential energy barrier, and a well for state B. In the overdamped (or "diffusive") regime, the transition rate from state A to B is related to the resonant frequency of the wells via where is the frequency of the well for state A, is the frequency of the barrier well, is the viscous damping, is the energy of the top of the barrier, is the energy of bottom of the well for state A, and is the temperature of the system times the Boltzmann constant. For general damping (overdamped or underdamped), there is a similar formula. Justification for the Eyring equation One of the most important features introduced by Eyring, Polanyi and Evans was the notion that activated complexes are in quasi-equilibrium with the reactants. The rate is then directly proportional to the concentration of these complexes multiplied by the frequency (kBT/h) with which they are converted into products. Below, a non-rigorous plausibility argument is given for the functional form of the Eyring equation. However, the key statistical mechanical factor kBT/h will not be justified, and the argument presented below does not constitute a true "derivation" of the Eyring equation. Quasi-equilibrium assumption Quasi-equilibrium is different from classical chemical equilibrium, but can be described using a similar thermodynamic treatment. Consider the reaction below {A} + {B} <=> {[AB]^\ddagger} -> {P} where complete equilibrium is achieved between all the species in the system including activated complexes, [AB]‡ . Using statistical mechanics, concentration of [AB]‡ can be calculated in terms of the concentration of A and B. TST assumes that even when the reactants and products are not in equilibrium with each other, the activated complexes are in quasi-equilibrium with the reactants. As illustrated in Figure 2, at any instant of time, there are a few activated complexes, and some were reactant molecules in the immediate past, which are designated [ABl]‡ (since they are moving from left to right). The remainder of them were product molecules in the immediate past ([ABr]‡). In TST, it is assumed that the flux of activated complexes in the two directions are independent of each other. That is, if all the product molecules were suddenly removed from the reaction system, the flow of [ABr]‡ stops, but there is still a flow from left to right. Hence, to be technically correct, the reactants are in equilibrium only with [ABl]‡, the activated complexes that were reactants in the immediate past. Plausibility argument The activated complexes do not follow a Boltzmann distribution of energies, but an "equilibrium constant" can still be derived from the distribution they do follow. The equilibrium constant K‡ for the quasi-equilibrium can be written as . So, the chemical activity of the transition state AB‡ is . Therefore, the rate equation for the production of product is , where the rate constant k is given by . Here, k‡ is directly proportional to the frequency of the vibrational mode responsible for converting the activated complex to the product; the frequency of this vibrational mode is . Every vibration does not necessarily lead to the formation of product, so a proportionality constant , referred to as the transmission coefficient, is introduced to account for this effect. So k‡ can be rewritten as . For the equilibrium constant K‡ , statistical mechanics leads to a temperature dependent expression given as (). Combining the new expressions for k‡ and K‡, a new rate constant expression can be written, which is given as . Since, by definition, ΔG‡ = ΔH‡ –TΔS‡, the rate constant expression can be expanded, to give an alternative form of the Eyring equation: . For correct dimensionality, the equation needs to have an extra factor of (c⊖)1–m for reactions that are not unimolecular: , where c⊖ is the standard concentration 1 mol⋅L−1 and m is the molecularity. Inferences from TST and relationship with Arrhenius theory The rate constant expression from transition state theory can be used to calculate the ΔG‡, ΔH‡, ΔS‡, and even ΔV‡ (the volume of activation) using experimental rate data. These so-called activation parameters give insight into the nature of a transition state, including energy content and degree of order, compared to the starting materials and has become a standard tool for elucidation of reaction mechanisms in physical organic chemistry. The free energy of activation, ΔG‡, is defined in transition state theory to be the energy such that holds. The parameters ΔH‡ and ΔS‡ can then be inferred by determining ΔG‡ = ΔH‡ – TΔS‡ at different temperatures. Because the functional form of the Eyring and Arrhenius equations are similar, it is tempting to relate the activation parameters with the activation energy and pre-exponential factors of the Arrhenius treatment. However, the Arrhenius equation was derived from experimental data and models the macroscopic rate using only two parameters, irrespective of the number of transition states in a mechanism. In contrast, activation parameters can be found for every transition state of a multistep mechanism, at least in principle. Thus, although the enthalpy of activation, ΔH‡, is often equated with Arrhenius's activation energy Ea, they are not equivalent. For a condensed-phase (e.g., solution-phase) or unimolecular gas-phase reaction step, Ea = ΔH‡ + RT. For other gas-phase reactions, Ea = ΔH‡ + (1 − Δn‡)RT, where Δn‡ is the change in the number of molecules on forming the transition state. (Thus, for a bimolecular gas-phase process, Ea = ΔH‡ + 2RT.) The entropy of activation, ΔS‡, gives the extent to which transition state (including any solvent molecules involved in or perturbed by the reaction) is more disordered compared to the starting materials. It offers a concrete interpretation of the pre-exponential factor A in the Arrhenius equation; for a unimolecular, single-step process, the rough equivalence A = (kBT/h) exp(1 + ΔS‡/R) (or A = (kBT/h) exp(2 + ΔS‡/R) for bimolecular gas-phase reactions) holds. For a unimolecular process, a negative value indicates a more ordered, rigid transition state than the ground state, while a positive value reflects a transition state with looser bonds and/or greater conformational freedom. It is important to note that, for reasons of dimensionality, reactions that are bimolecular or higher have ΔS‡ values that depend on the standard state chosen (standard concentration, in particular). For most recent publications, 1 mol L−1 or 1 molar is chosen. Since this choice is a human construct, based on our definitions of units for molar quantity and volume, the magnitude and sign of ΔS‡ for a single reaction is meaningless by itself; only comparisons of the value with that of a reference reaction of "known" (or assumed) mechanism, made at the same standard state, is valid. The volume of activation is found by taking the partial derivative of ΔG‡ with respect to pressure (holding temperature constant): . It gives information regarding the size, and hence, degree of bonding at the transition state. An associative mechanism will likely have a negative volume of activation, while a dissociative mechanism will likely have a positive value. Given the relationship between equilibrium constant and the forward and reverse rate constants, , the Eyring equation implies that . Another implication of TST is the Curtin–Hammett principle: the product ratio of a kinetically-controlled reaction from R to two products A and B will reflect the difference in the energies of the respective transition states leading to product, assuming there is a single transition state to each one: (). (In the expression for ΔΔG‡ above, there is an extra term if A and B are formed from two different species SA and SB in equilibrium.) For a thermodynamically-controlled reaction, every difference of RT ln 10 ≈ (1.987 × 10−3 kcal/mol K)(298 K)(2.303) ≈ 1.36 kcal/mol in the free energies of products A and B results in a factor of 10 in selectivity at room temperature (298 K), a principle known as the "1.36 rule": (). Analogously, every 1.36 kcal/mol difference in the free energy of activation results in a factor of 10 in selectivity for a kinetically-controlled process at room temperature: (). Using the Eyring equation, there is a straightforward relationship between ΔG‡, first-order rate constants, and reaction half-life at a given temperature. At 298 K, a reaction with ΔG‡ = 23 kcal/mol has a rate constant of k ≈ 8.4 × 10−5 s−1 and a half life of t1/2 ≈ 2.3 hours, figures that are often rounded to k ~ 10−4 s−1 and t1/2 ~ 2 h. Thus, a free energy of activation of this magnitude corresponds to a typical reaction that proceeds to completion overnight at room temperature. For comparison, the cyclohexane chair flip has a ΔG‡ of about 11 kcal/mol with k ~ 105 s−1, making it a dynamic process that takes place rapidly (faster than the NMR timescale) at room temperature. At the other end of the scale, the cis/trans isomerization of 2-butene has a ΔG‡ of about 60 kcal/mol, corresponding to k ~ 10−31 s−1 at 298 K. This is a negligible rate: the half-life is 12 orders of magnitude longer than the age of the universe. Limitations In general, TST has provided researchers with a conceptual foundation for understanding how chemical reactions take place. Even though the theory is widely applicable, it does have limitations. For example, when applied to each elementary step of a multi-step reaction, the theory assumes that each intermediate is long-lived enough to reach a Boltzmann distribution of energies before continuing to the next step. When the intermediates are very short-lived, TST fails. In such cases, the momentum of the reaction trajectory from the reactants to the intermediate can carry forward to affect product selectivity. An example of such a reaction is the ring closure of cyclopentane biradicals generated from the gas-phase thermal decomposition of 2,3-diazabicyclo[2.2.1]hept-2-ene. Transition state theory is also based on the assumption that atomic nuclei behave according to classical mechanics. It is assumed that unless atoms or molecules collide with enough energy to form the transition structure, then the reaction does not occur. However, according to quantum mechanics, for any barrier with a finite amount of energy, there is a possibility that particles can still tunnel across the barrier. With respect to chemical reactions this means that there is a chance that molecules will react, even if they do not collide with enough energy to overcome the energy barrier. While this effect is negligible for reactions with large activation energies, it becomes an important phenomenon for reactions with relatively low energy barriers, since the tunneling probability increases with decreasing barrier height. Transition state theory fails for some reactions at high temperature. The theory assumes the reaction system will pass over the lowest energy saddle point on the potential energy surface. While this description is consistent for reactions occurring at relatively low temperatures, at high temperatures, molecules populate higher energy vibrational modes; their motion becomes more complex and collisions may lead to transition states far away from the lowest energy saddle point. This deviation from transition state theory is observed even in the simple exchange reaction between diatomic hydrogen and a hydrogen radical. Given these limitations, several alternatives to transition state theory have been proposed. A brief discussion of these theories follows. Generalized transition state theory Any form of TST, such as microcanonical variational TST, canonical variational TST, and improved canonical variational TST, in which the transition state is not necessarily located at the saddle point, is referred to as generalized transition state theory. Microcanonical variational TST A fundamental flaw of transition state theory is that it counts any crossing of the transition state as a reaction from reactants to products or vice versa. In reality, a molecule may cross this "dividing surface" and turn around, or cross multiple times and only truly react once. As such, unadjusted TST is said to provide an upper bound for the rate coefficients. To correct for this, variational transition state theory varies the location of the dividing surface that defines a successful reaction in order to minimize the rate for each fixed energy. The rate expressions obtained in this microcanonical treatment can be integrated over the energy, taking into account the statistical distribution over energy states, so as to give the canonical, or thermal rates. Canonical variational TST A development of transition state theory in which the position of the dividing surface is varied so as to minimize the rate constant at a given temperature. Improved canonical variational TST A modification of canonical variational transition state theory in which, for energies below the threshold energy, the position of the dividing surface is taken to be that of the microcanonical threshold energy. This forces the contributions to rate constants to be zero if they are below the threshold energy. A compromise dividing surface is then chosen so as to minimize the contributions to the rate constant made by reactants having higher energies. Nonadiabatic TST An expansion of TST to the reactions when two spin-states are involved simultaneously is called nonadiabatic transition state theory (NA-TST). Semiclassical TST Using vibrational perturbation theory, effects such as tunnelling and variational effects can be accounted for within the SCTST formalism. Applications Enzymatic reactions Enzymes catalyze chemical reactions at rates that are astounding relative to uncatalyzed chemistry at the same reaction conditions. Each catalytic event requires a minimum of three or often more steps, all of which occur within the few milliseconds that characterize typical enzymatic reactions. According to transition state theory, the smallest fraction of the catalytic cycle is spent in the most important step, that of the transition state. The original proposals of absolute reaction rate theory for chemical reactions defined the transition state as a distinct species in the reaction coordinate that determined the absolute reaction rate. Soon thereafter, Linus Pauling proposed that the powerful catalytic action of enzymes could be explained by specific tight binding to the transition state species Because reaction rate is proportional to the fraction of the reactant in the transition state complex, the enzyme was proposed to increase the concentration of the reactive species. This proposal was formalized by Wolfenden and coworkers at University of North Carolina at Chapel Hill, who hypothesized that the rate increase imposed by enzymes is proportional to the affinity of the enzyme for the transition state structure relative to the Michaelis complex. Because enzymes typically increase the non-catalyzed reaction rate by factors of 106-1026, and Michaelis complexes often have dissociation constants in the range of 10−3-10−6 M, it is proposed that transition state complexes are bound with dissociation constants in the range of 10−14 -10−23 M. As substrate progresses from the Michaelis complex to product, chemistry occurs by enzyme-induced changes in electron distribution in the substrate. Enzymes alter the electronic structure by protonation, proton abstraction, electron transfer, geometric distortion, hydrophobic partitioning, and interaction with Lewis acids and bases. Analogs that resemble the transition state structures should therefore provide the most powerful noncovalent inhibitors known. All chemical transformations pass through an unstable structure called the transition state, which is poised between the chemical structures of the substrates and products. The transition states for chemical reactions are proposed to have lifetimes near 10−13 seconds, on the order of the time of a single bond vibration. No physical or spectroscopic method is available to directly observe the structure of the transition state for enzymatic reactions, yet transition state structure is central to understanding enzyme catalysis since enzymes work by lowering the activation energy of a chemical transformation. It is now accepted that enzymes function to stabilize transition states lying between reactants and products, and that they would therefore be expected to bind strongly any inhibitor that closely resembles such a transition state. Substrates and products often participate in several enzyme catalyzed reactions, whereas the transition state tends to be characteristic of one particular enzyme, so that such an inhibitor tends to be specific for that particular enzyme. The identification of numerous transition state inhibitors supports the transition state stabilization hypothesis for enzymatic catalysis. Currently there is a large number of enzymes known to interact with transition state analogs, most of which have been designed with the intention of inhibiting the target enzyme. Examples include HIV-1 protease, racemases, β-lactamases, metalloproteinases, cyclooxygenases and many others. Adsorption on surfaces and reactions on surfaces Desorption as well as reactions on surfaces are straightforward to describe with transition state theory. Analysis of adsorption to a surface from a liquid phase can present a challenge due to lack of ability to assess the concentration of the solute near the surface. When full details are not available, it has been proposed that reacting species' concentrations should be normalized to the concentration of active surface sites, an approximation called the surface reactant equi-density approximation (SREA).	Physical sciences	Kinetics	Chemistry
9221221	https://en.wikipedia.org/wiki/Ceramic%20capacitor	Ceramic capacitor	A ceramic capacitor is a fixed-value capacitor where the ceramic material acts as the dielectric. It is constructed of two or more alternating layers of ceramic and a metal layer acting as the electrodes. The composition of the ceramic material defines the electrical behavior and therefore applications. Ceramic capacitors are divided into two application classes: Class 1 ceramic capacitors offer high stability and low losses for resonant circuit applications. Class 2 ceramic capacitors offer high volumetric efficiency for buffer, by-pass, and coupling applications. Ceramic capacitors, especially multilayer ceramic capacitors (MLCCs), are the most produced and used capacitors in electronic equipment that incorporate approximately one trillion (1012) pieces per year. Ceramic capacitors of special shapes and styles are used as capacitors for RFI/EMI suppression, as feed-through capacitors and in larger dimensions as power capacitors for transmitters. History Since the beginning of the study of electricity non-conductive materials such as glass, porcelain, paper and mica have been used as insulators. These materials some decades later were also well-suited for further use as the dielectric for the first capacitors. Even in the early years of Marconi's wireless transmitting apparatus, porcelain capacitors were used for high voltage and high frequency application in the transmitters. On the receiver side, the smaller mica capacitors were used for resonant circuits. Mica dielectric capacitors were invented in 1909 by William Dubilier. Prior to World War II, mica was the most common dielectric for capacitors in the United States. Mica is a natural material and not available in unlimited quantities. So in the mid-1920s the deficiency of mica in Germany and the experience in porcelain—a special class of ceramic—led in Germany to the first capacitors using ceramic as dielectric, founding a new family of ceramic capacitors. Paraelectric titanium dioxide (rutile) was used as the first ceramic dielectric because it had a linear temperature dependence of capacitance for temperature compensation of resonant circuits and can replace mica capacitors. In 1926 these ceramic capacitors were produced in small quantities with increasing quantities in the 1940s. The style of these early ceramics was a disc with metallization on both sides contacted with tinned wires. This style predates the transistor and was used extensively in vacuum-tube equipment (e.g., radio receivers) from about 1930 through the 1950s. But this paraelectric dielectric had relatively low permittivity so that only small capacitance values could be realized. The expanding market of radios in the 1930s and 1940s create a demand for higher capacitance values but below electrolytic capacitors for HF decoupling applications. Discovered in 1921, the ferroelectric ceramic material barium titanate with a permittivity in the range of 1,000, about ten times greater than titanium dioxide or mica, began to play a much larger role in electronic applications. The higher permittivity resulted in much higher capacitance values, but this was coupled with relatively unstable electrical parameters. Therefore, these ceramic capacitors only could replace the commonly used mica capacitors for applications where stability was less important. Smaller dimensions, as compared to the mica capacitors, lower production costs and independence from mica availability accelerated their acceptance. The fast-growing broadcasting industry after the Second World War drove deeper understanding of the crystallography, phase transitions and the chemical and mechanical optimization of the ceramic materials. Through the complex mixture of different basic materials, the electrical properties of ceramic capacitors can be precisely adjusted. To distinguish the electrical properties of ceramic capacitors, standardization defined several different application classes (Class 1, Class 2, Class 3). It is remarkable that the separate development during the War and the time afterwards in the US and the European market had led to different definitions of these classes (EIA vs IEC), and only recently (since 2010) has a worldwide harmonization to the IEC standardization taken place. The typical style for ceramic capacitors beneath the disc (at that time called condensers) in radio applications at the time after the War from the 1950s through the 1970s was a ceramic tube covered with tin or silver on both the inside and outside surface. It included relatively long terminals forming, together with resistors and other components, a tangle of open circuit wiring. The easy-to-mold ceramic material facilitated the development of special and large styles of ceramic capacitors for high-voltage, high-frequency (RF) and power applications. With the development of semiconductor technology in the 1950s, barrier layer capacitors, or IEC class 3/EIA class IV capacitors, were developed using doped ferroelectric ceramics. Because this doped material was not suitable to produce multilayers, they were replaced decades later by Y5V class 2 capacitors. The early style of the ceramic disc capacitor could be more cheaply produced than the common ceramic tube capacitors in the 1950s and 1970s. An American company in the midst of the Apollo program, launched in 1961, pioneered the stacking of multiple discs to create a monolithic block. This "multi-layer ceramic capacitor" (MLCC) was compact and offered high-capacitance capacitors. The production of these capacitors using the tape casting and ceramic-electrode cofiring processes was a great manufacturing challenge. MLCCs expanded the range of applications to those requiring larger capacitance values in smaller cases. These ceramic chip capacitors were the driving force behind the conversion of electronic devices from through-hole mounting to surface-mount technology in the 1980s. Polarized electrolytic capacitors could be replaced by non-polarized ceramic capacitors, simplifying the mounting. In 1993, TDK Corporation succeeded in displacing palladium bearing electrodes with much cheaper nickel electrodes, significantly reducing production costs and enabling mass production of MLCCs. , more than 1012 MLCCs are manufactured each year. Along with the style of ceramic chip capacitors, ceramic disc capacitors are often used as safety capacitors in electromagnetic interference suppression applications. Besides these, large ceramic power capacitors for high voltage or high frequency transmitter applications are also to be found. New developments in ceramic materials have been made with anti-ferroelectric ceramics. This material has a nonlinear antiferroelectric/ferroelectric phase change that allows increased energy storage with higher volumetric efficiency. They are used for energy storage (for example, in detonators). Application classes, definitions The different ceramic materials used for ceramic capacitors, paraelectric or ferroelectric ceramics, influences the electrical characteristics of the capacitors. Using mixtures of paraelectric substances based on titanium dioxide results in very stable and linear behavior of the capacitance value within a specified temperature range and low losses at high frequencies. But these mixtures have a relatively low permittivity so that the capacitance values of these capacitors are relatively small. Higher capacitance values for ceramic capacitors can be attained by using mixtures of ferroelectric materials like barium titanate together with specific oxides. These dielectric materials have much higher permittivities, but at the same time their capacitance value are more or less nonlinear over the temperature range, and losses at high frequencies are much higher. These different electrical characteristics of ceramic capacitors requires to group them into "application classes". The definition of the application classes comes from the standardization. As of 2013, two sets of standards were in use, one from International Electrotechnical Commission (IEC) and the other from the now-defunct Electronic Industries Alliance (EIA). The definitions of the application classes given in the two standards are different. The following table shows the different definitions of the application classes for ceramic capacitors: Manufacturers, especially in the US, preferred Electronic Industries Alliance (EIA) standards. In many parts very similar to the IEC standard, the EIA RS-198 defines four application classes for ceramic capacitors. The different class numbers within both standards are the reason for a lot of misunderstandings interpreting the class descriptions in the datasheets of many manufacturers. The EIA ceased operations on February 11, 2011, but the former sectors continue to serve international standardization organizations. In the following, the definitions of the IEC standard will be preferred and in important cases compared with the definitions of the EIA standard. Class 1 ceramic capacitors Class 1 ceramic capacitors are accurate, temperature-compensating capacitors. They offer the most stable voltage, temperature, and to some extent, frequency. They have the lowest losses and therefore are especially suited for resonant circuit applications where stability is essential or where a precisely defined temperature coefficient is required, for example in compensating temperature effects for a circuit. The basic materials of class 1 ceramic capacitors are composed of a mixture of finely ground granules of paraelectric materials such as titanium dioxide (), modified by additives of zinc, zirconium, niobium, magnesium, tantalum, cobalt and strontium, which are necessary to achieve the capacitor's desired linear characteristics. The general capacitance temperature behavior of class 1 capacitors depends on the basic paraelectric material, for example . The additives of the chemical composition are used to adjust precisely the desired temperature characteristic. Class 1 ceramic capacitors have the lowest volumetric efficiency among ceramic capacitors. This is the result of the relatively low permittivity (6 to 200) of paraelectric materials. Therefore, class 1 capacitors have capacitance values in the lower range. Class 1 capacitors have a temperature coefficient that is typically fairly linear with temperature. These capacitors have very low electrical losses with a dissipation factor of approximately 0.15%. They undergo no significant aging processes and the capacitance value is nearly independent of the applied voltage. These characteristics allow applications for high Q filters, in resonant circuits and oscillators (for example, in phase-locked loop circuits). The EIA RS-198 standard codes ceramic class 1 capacitors with a three character code that indicates temperature coefficient. The first letter gives the significant figure of the change in capacitance over temperature (temperature coefficient α) in ppm/K. The second character gives the multiplier of the temperature coefficient. The third letter gives the maximum tolerance from that in ppm/K. All ratings are from 25 to 85 °C: In addition to the EIA code, the temperature coefficient of the capacitance dependence of class 1 ceramic capacitors is commonly expressed in ceramic names like "NP0", "N220" etc. These names include the temperature coefficient (α). In the IEC/EN 60384-8/21 standard, the temperature coefficient and tolerance are replaced by a two digit letter code (see table) in which the corresponding EIA code is added. For instance, an "NP0" capacitor with EIA code "C0G" will have 0 drift, with a tolerance of ±30 ppm/K, while an "N1500" with the code "P3K" will have −1500 ppm/K drift, with a maximum tolerance of ±250 ppm/K. Note that the IEC and EIA capacitor codes are industry capacitor codes and not the same as military capacitor codes. Class 1 capacitors include capacitors with different temperature coefficients α. Especially, NP0/CG/C0G capacitors with an α ±0•10−6 /K and an α tolerance of 30 ppm are technically of great interest. These capacitors have a capacitance variation dC/C of ±0.54% within the temperature range −55 to +125 °C. This enables accurate frequency response over a wide temperature range (in, for example, resonant circuits). The other materials with their special temperature behavior are used to compensate a counter temperature run of parallel connected components like coils in oscillator circuits. Class 1 capacitors exhibit very small tolerances of the rated capacitance. Class 2 ceramic capacitors Class 2 ceramic capacitors have a dielectric with a high permittivity and therefore a better volumetric efficiency than class 1 capacitors, but lower accuracy and stability. The ceramic dielectric is characterized by a nonlinear change of capacitance over the temperature range. The capacitance value also depends on the applied voltage. They are suitable for bypass, coupling and decoupling applications or for frequency discriminating circuits where low losses and high stability of capacitance are less important. They typically exhibit microphony. Class 2 capacitors are made of ferroelectric materials such as barium titanate () and suitable additives such as aluminium silicate, magnesium silicate and aluminium oxide. These ceramics have very high permittivity (200 to 14,000), allowing an extreme electric field and therefore capacitance within relatively small packages — class 2 capacitors are significantly smaller than comparable class 1 capacitors. However, the permittivity is nonlinear with respect to field strength, meaning the capacitance varies significantly as the voltage across the terminals increases. Class 2 capacitors also exhibit poor temperature stability and age over time. Due to these traits, class 2 capacitors are typically used in applications where only a minimum value of capacitance (as opposed to an accurate value) is required, such as the buffering/filtering of inputs and outputs of power supplies, and the coupling of electric signals. Class 2 capacitors are labeled according to the change in capacitance over the temperature range. The most widely used classification is based on the EIA RS-198 standard and uses a three-digit code. The first character, a letter, denotes the coldest operating temperature; the second character, a numeral, denotes the hottest temperature; and the third character, another letter, denotes the maximum allowed capacitance change over the capacitor's entire specified temperature range: For instance, a Z5U capacitor will operate from +10 °C to +85 °C with a capacitance change of at most +22% to −56%. An X7R capacitor will operate from −55 °C to +125 °C with a capacitance change of at most ±15%. Some commonly used class 2 ceramic capacitor materials are listed below: X8R (−55/+150, ΔC/C0 = ±15%), X7R (−55/+125 °C, ΔC/C0 = ±15%), X6R (−55/+105 °C, ΔC/C0 = ±15%), X5R (−55/+85 °C, ΔC/C0 = ±15%), X7S (−55/+125, ΔC/C0 = ±22%), Z5U (+10/+85 °C, ΔC/C0 = +22/−56%), Y5V (−30/+85 °C, ΔC/C0 = +22/−82%), The IEC/EN 60384 -9/22 standard uses another two-digit-code. In most cases it is possible to translate the EIA code into the IEC/EN code. Slight translation errors occur, but normally are tolerable. X7R correlates with 2X1 Z5U correlates with 2E6 Y5V similar to 2F4, aberration: ΔC/C0 = +30/−80% instead of +30/−82% X7S similar to 2C1, aberration: ΔC/C0 = ±20% instead of ±22% X8R no IEC/EN code available Because class 2 ceramic capacitors have lower capacitance accuracy and stability, they require higher tolerance. For military types the class 2 dielectrics specify temperature characteristic (TC) but not temperature-voltage characteristic (TVC). Similar to X7R, military type BX cannot vary more than 15% over temperature, and in addition, must remain within +15%/-25 % at maximum rated voltage. Type BR has a TVC limit of +15%/-40%. Class 3 ceramic capacitors Class 3 barrier layer or semiconductive ceramic capacitors have very high permittivity, up to 50,000 and therefore a better volumetric efficiency than class 2 capacitors. However, these capacitors have worse electrical characteristics, including lower accuracy and stability. The dielectric is characterized by very high nonlinear change of capacitance over the temperature range. The capacitance value additionally depends on the voltage applied. As well, they have very high losses and age over time. Barrier layer ceramic capacitors are made of doped ferroelectric materials such as barium titanate (). As this ceramic technology improved in the mid-1980s, barrier layer capacitors became available in values of up to 100 μF, and at that time it seemed that they could substitute for smaller electrolytic capacitors. Because it is not possible to build multilayer capacitors with this material, only leaded single layer types are offered in the market. Due to advancements in multilayer ceramic capacitors enabling superior performance in a smaller package, barrier layer capacitors as a technology are now considered obsolete and no longer standardized by the IEC. Construction and styles Ceramic capacitors are composed of a mixture of finely ground granules of paraelectric or ferroelectric materials, appropriately mixed with other materials to achieve the desired characteristics. From these powder mixtures, the ceramic is sintered at high temperatures. The ceramic forms the dielectric and serves as a carrier for the metallic electrodes. The minimum thickness of the dielectric layer, which today (2013) for low voltage capacitors is in the size range of 0.5 micrometers is limited downwards by the grain size of the ceramic powder. The thickness of the dielectric for capacitors with higher voltages is determined by the dielectric strength of the desired capacitor. The electrodes of the capacitor are deposited on the ceramic layer by metallization. For MLCCs alternating metallized ceramic layers are stacked one above the other. The outstanding metallization of the electrodes at both sides of the body are connected with the contacting terminal. A lacquer or ceramic coating protects the capacitor against moisture and other ambient influences. Ceramic capacitors come in various shapes and styles. Some of the most common are: Multilayer ceramic chip capacitor (MLCC), rectangular block, for surface mounting Ceramic disc capacitor, single layer disc, resin coated, with through-hole leads Feedthrough ceramic capacitor, used for bypass purposes in high-frequency circuits. Tube shape, inner metallization contacted with a lead, outer metallization for soldering Ceramic power capacitors, larger ceramic bodies in different shapes for high voltage applications Multi-layer ceramic capacitors (MLCC) Manufacturing An MLCC can be thought of as consisting of many single-layer capacitors stacked together into a single package. The starting material for all MLCC chips is a mixture of finely ground granules of paraelectric or ferroelectric raw materials, modified by accurately determined additives. The composition of the mixture and the size of the powder particles, as small as 10 nm, reflect the manufacturer's expertise. A thin ceramic foil is cast from a suspension of the powder with a suitable binder. Rolls of foil are cut into equal-sized sheets, which are screen printed with a metal paste layer, which will become the electrodes. In an automated process, these sheets are stacked in the required number of layers and solidified by pressure. Besides the relative permittivity, the size and number of layers determines the later capacitance value. The electrodes are stacked in an alternating arrangement slightly offset from the adjoining layers so that they each can later be connected on the offset side, one left, one right. The layered stack is pressed and then cut into individual components. High mechanical precision is required, for example, to produce a 500 or more layer stack of size "0201" (0.5 mm × 0.3 mm). After cutting, the binder is burnt out of the stack. This is followed by sintering at temperatures between , producing the final, mainly crystalline, structure. This burning process creates the desired dielectric properties. Burning is followed by cleaning and then metallization of both end surfaces. Through the metallization, the ends and the inner electrodes are connected in parallel and the capacitor gets its terminals. Finally, each capacitor is electrically tested to ensure functionality and adequate performance, and packaged in a tape reel. Miniaturizing The capacitance formula (C) of a MLCC capacitor is based on the formula for a plate capacitor enhanced with the number of layers: where ε stands for dielectric permittivity; A for electrode surface area; n for the number of layers; and d for the distance between the electrodes. A thinner dielectric or a larger electrode area each increase the capacitance value, as will a dielectric material of higher permittivity. With the progressive miniaturization of digital electronics in recent decades, the components on the periphery of the integrated logic circuits have been scaled down as well. Shrinking an MLCC involves reducing the dielectric thickness and increasing the number of layers. Both options require huge efforts and are connected with a lot of expertise. In 1995 the minimum thickness of the dielectric was 4 μm. By 2005 some manufacturers produced MLCC chips with layer thicknesses of 1 μm. , the minimum thickness is about 0.5 μm. The field strength in the dielectric increased to 35 V/μm. The size reduction of these capacitors is achieved reducing powder grain size, the assumption to make the ceramic layers thinner. In addition, the manufacturing process became more precisely controlled, so that more and more layers can be stacked. Between 1995 and 2005, the capacitance of a Y5V MLCC capacitor of size 1206 was increased from 4.7 μF to 100 μF. Meanwhile, (2013) a lot of producers can deliver class 2 MLCC capacitors with a capacitance value of 100 μF in the chip-size 0805. MLCC case sizes MLCCs don't have leads, and as a result they are usually smaller than their counterparts with leads. They don't require through-hole access in a PCB to mount and are designed to be handled by machines rather than by humans. As a result, surface-mount components like MLCCs are typically cheaper. MLCCs are manufactured in standardized shapes and sizes for comparable handling. Because the early standardization was dominated by American EIA standards the dimensions of the MLCC chips were standardized by EIA in units of inches. A rectangular chip with the dimensions of 0.06-inch length and 0.03-inch width is coded as "0603". This code is international and in common use. JEDEC (IEC/EN), devised a second, metric code. The EIA code and the metric equivalent of the common sizes of multilayer ceramic chip capacitors, and the dimensions in mm are shown in the following table. Missing from the table is the measure of the height "H". This is generally not listed, because the height of MLCC chips depends on the number of layers and thus on the capacitance. Normally, however, the height H does not exceed the width W. NME and BME metallization Originally, MLCC electrodes were constructed out of noble metals such as silver and palladium which can withstand high sintering temperatures of without readily oxidizing. These noble metal electrode (NME) capacitors offered very good electrical properties. However, a surge in prices of noble metals in the late 1990s greatly increased manufacturing costs; these pressures resulted in the development of capacitors that used cheaper metals like copper and nickel. These base metal electrode (BME) capacitors possessed poorer electrical characteristics; exhibiting greater shrinkage of capacitance at higher voltages and increased loss factor. The disadvantages of BME were deemed acceptable for class 2 capacitors, which are primarily used in accuracy-insensitive, low-cost applications such as power supplies. NME still sees use in class 1 capacitors where conformance to specifications are critical and cost is less of a concern. MLCC capacitance ranges Capacitance of MLCC chips depends on the dielectric, the size and the required voltage (rated voltage). Capacitance values start at about 1pF. The maximum capacitance value is determined by the production technique. For X7R that is 47 μF, for Y5V: 100 μF. The picture right shows the maximum capacitance for class 1 and class 2 multilayer ceramic chip capacitors. The following two tables, for ceramics NP0/C0G and X7R each, list for each common case size the maximum available capacitance value and rated voltage of the leading manufacturers Murata, TDK, KEMET, AVX. (Status April 2017) Low-ESL styles In the region of its resonance frequency, a capacitor has the best decoupling properties for noise or electromagnetic interference. The resonance frequency of a capacitor is determined by the inductance of the component. The inductive parts of a capacitor are summarized in the equivalent series inductance, or ESL. (Note that L is the electrical symbol for inductance.) The smaller the inductance, the higher the resonance frequency. Because, especially in digital signal processing, switching frequencies have continued to rise, the demand for high frequency decoupling or filter capacitors increases. With a simple design change the ESL of an MLCC chip can be reduced. Therefore, the stacked electrodes are connected on the longitudinal side with the connecting terminations. This reduces the distance that the charge carriers flow over the electrodes, which reduces inductance of the component. For example, an 0.1 μF X7R MLCC in a 0805 package resonates at 16 MHz. The same capacitor with leads on its long sides (i.e. an 0508) has a resonance frequency of 22 MHz. Another possibility is to form the device as an array of capacitors. Here, several individual capacitors are built in a common housing. Connecting them in parallel, the resulting ESL as well as ESR values of the components are reduced. X2Y decoupling capacitor A standard multi-layer ceramic capacitor has many opposing electrode layers stacked inside connected with two outer terminations. The X2Y ceramic chip capacitor however is a 4 terminal chip device. It is constructed like a standard two-terminal MLCC out of the stacked ceramic layers with an additional third set of shield electrodes incorporated in the chip. These shield electrodes surround each existing electrode within the stack of the capacitor plates and are low ohmic contacted with two additional side terminations across to the capacitor terminations. The X2Y construction results in a three-node capacitive circuit that provides simultaneous line-to-line and line-to-ground filtering. Capable of replacing 2 or more conventional devices, the X2Y ceramic capacitors are ideal for high frequency filtering or noise suppression of supply voltages in digital circuits, and can prove invaluable in meeting stringent EMC demands in dc motors, in automotive, audio, sensor and other applications. The X2Y footprint results in lower mounted inductance. This is particularly of interest for use in high-speed digital circuits with clock rates of several 100 MHz and upwards. There the decoupling of the individual supply voltages on the circuit board is difficult to realize due to parasitic inductances of the supply lines. A standard solution with conventional ceramic capacitors requires the parallel use of many conventional MLCC chips with different capacitance values. Here X2Y capacitors are able to replace up to five equal-sized ceramic capacitors on the PCB. However, this particular type of ceramic capacitor is patented, so these components are still comparatively expensive. An alternative to X2Y capacitors may be a three-terminal capacitor. Mechanical susceptibility Ceramics are brittle, and MLCC chips surface-mount soldered to a circuit board are often vulnerable to cracking from thermal expansion or mechanical stresses like depanelization, more so than leaded through-hole components. The cracks can come from automated machine assembly line, or from high current in the circuit. Vibration and shock forces on the circuit board are more or less transmitted undampened to the MLCC and its solder joints; excessive force may cause the capacitor to crack (flex crack). Excess solder in the joints are undesirable as they may magnify the forces that the capacitor is subject to. The capability of MLCC chips to withstand mechanical stress is tested by a so-called substrate bending test, where a PCB with a soldered MLCC is bent by a punch by 1 to 3 mm. Failure occurs if the MLCC becomes a short-circuit or significantly changes in capacitance. Bending strengths of MLCC chips differ by the ceramic material, the size of the chip, and the physical construction of the capacitors. Without special mitigation, NP0/C0G class 1 ceramic MLCC chips reach a typical bending strength of 2 mm while larger types of X7R, Y5V class 2 ceramic chips achieved only a bending strength of approximately 1 mm. Smaller chips, such as the size of 0402, reached in all types of ceramics larger bending strength values. With special design features, particularly at the electrodes and terminations, the bending strength can be improved. For example, an internal short circuit arises by the contact of two electrodes with opposite polarity, which will be produced at the break of the ceramic in the region of the terminations. This can be prevented when the overlap surfaces of the electrodes are reduced. This is achieved e.g. by an "Open Mode Design" (OMD). Here a break in the region of the terminations only reduce the capacitance value a little bit (AVX, KEMET). With a similar construction called "Floating Electrode Design" (FED) or "Multi-layer Serial Capacitors" (MLSC), also, only capacitance reduction results if parts of the capacitor body break. This construction works with floating electrodes without any conductive connection to the termination. A break doesn't lead to a short, only to capacitance reduction. However, both structures lead to larger designs with respect to a standard MLCC version with the same capacitance value. The same volume with respect to standard MLCCs is achieved by the introduction of a flexible intermediate layer of a conductive polymer between the electrodes and the termination called "Flexible Terminations" (FT-Cap) or "Soft Terminations". In this construction, the rigid metallic soldering connection can move against the flexible polymer layer, and thus can absorb the bending forces, without resulting in a break in the ceramic. Some automotive capacitors are specified to adhere to AEC-Q200 and/or VW 80808. RFI/EMI suppression with X- and Y capacitors Suppression capacitors are effective interference reduction components because their electrical impedance decreases with increasing frequency, such that at higher frequencies they appear as short circuits to high-frequency electrical noise and transients between the lines, or to ground. They therefore prevent equipment and machinery (including motors, inverters, and electronic ballasts, as well as solid-state relay snubbers and spark quenchers) from sending and receiving electromagnetic and radio frequency interference as well as transients in across-the-line (X capacitors) and line-to-ground (Y capacitors) connections. X capacitors effectively absorb symmetrical, balanced, or differential interference. Y capacitors are connected in a line bypass between a line phase and a point of zero potential, to absorb asymmetrical, unbalanced, or common-mode interference. EMI/RFI suppression capacitors are designed so that any remaining interference or electrical noise does not exceed the limits of EMC directive EN 50081. Suppression components are connected directly to mains voltage for 10 to 20 years or more and are therefore exposed to potentially damaging overvoltages and transients. For this reason, suppression capacitors must comply with the safety and non-flammability requirements of international safety standards such as Europe: EN 60384-14, USA: UL 1414, UL 1283 Canada: CSA C22.2, No.1, CSA C22.2, No.8 China: CQC (GB/T 14472-1998) RFI capacitors that fulfill all specified requirements are imprinted with the certification mark of various national safety standards agencies. For power line applications, special requirements are placed on the non-flammability of the coating and the epoxy resin impregnating or coating the capacitor body. To receive safety approvals, X and Y powerline-rated capacitors are destructively tested to the point of failure. Even when exposed to large overvoltage surges, these safety-rated capacitors must fail in a fail-safe manner that does not endanger personnel or property. most ceramic capacitors used for EMI/RFI suppression were leaded ones for through-hole mounting on a PCB, the surface-mount technique is becoming more and more important. For this reason, in recent years a lot of MLCC chips for EMI/RFI suppression from different manufacturers have received approvals and fulfill all requirements given in the applicable standards. Ceramic power capacitors Although the materials used for large power ceramic capacitors mostly are very similar to those used for smaller ones, ceramic capacitors with high to very high power or voltage ratings for applications in power systems, transmitters and electrical installations are often classified separately, for historical reasons. The standardization of ceramic capacitors for lower power is oriented toward electrical and mechanical parameters as components for use in electronic equipment. The standardization of power capacitors, contrary to that, is strongly focused on protecting personnel and equipment, given by the local regulating authority. As modern electronic equipment gained the ability to handle power levels that were previously the exclusive domain of "electrical power" components, the distinction between the "electronic" and "electrical" power ratings has become less distinct. In the past, the boundary between these two families was approximately at a reactive power of 200 volt-amps, but modern power electronics can handle increasing amounts of power. Power ceramic capacitors are mostly specified for much higher than 200 volt-amps. The great plasticity of ceramic raw material and the high dielectric strength of ceramics deliver solutions for many applications and are the reasons for the enormous diversity of styles within the family of power ceramic capacitors. These power capacitors have been on the market for decades. They are produced according to the requirements as class 1 power ceramic capacitors with high stability and low losses or class 2 power ceramic capacitors with high volumetric efficiency. Class 1 power ceramic capacitors are used for resonant circuit application in transmitter stations. Class 2 power ceramic capacitors are used for circuit breakers, for power distribution lines, for high voltage power supplies in laser-applications, for induction furnaces and in voltage-doubling circuits. Power ceramic capacitors can be supplied with high rated voltages in the range of 2 kV up to 100 kV. The dimensions of these power ceramic capacitors can be very large. At high power applications the losses of these capacitors can generate a lot of heat. For this reason some special styles of power ceramic capacitors have pipes for water-cooling. Electrical characteristics Series-equivalent circuit All electrical characteristics of ceramic capacitors can be defined and specified by a series equivalent circuit composed out of an idealized capacitance and additional electrical components, which model all losses and inductive parameters of a capacitor. In this series-equivalent circuit the electrical characteristics of a capacitors is defined by C, the capacitance of the capacitor, Rinsul, the insulation resistance of the dielectric, not to be confused with the insulation of the housing RESR, the equivalent series resistance, which summarizes all ohmic losses of the capacitor, usually abbreviated as "ESR". LESL, the equivalent series inductance, which is the effective self-inductance of the capacitor, usually abbreviated as "ESL". The use of a series equivalent circuit instead of a parallel equivalent circuit is defined in IEC/EN 60384-1. Capacitance standard values and tolerances The "rated capacitance" CR or "nominal capacitance" CN is the value for which the capacitor has been designed. The actual capacitance depends on the measuring frequency and the ambient temperature. Standardized conditions for capacitors are a low-voltage AC measuring method at a temperature of 20 °C with frequencies of Class 1 ceramic capacitors CR ≤ 100 pF at 1 MHz, measuring voltage 5 V CR > 100 pF at 1 kHz, measuring voltage 5 V Class 2 ceramic capacitors CR ≤ 100 pF at 1 MHz, measuring voltage 1 V 100 pF < CR ≤ 10 μF at 1 kHz, measuring voltage 1 V CR > 10 μF at 100/120 Hz, measuring voltage 0.5 V Capacitors are available in different, geometrically increasing preferred values as specified in the E series standards specified in IEC/EN 60063. According to the number of values per decade, these were called the E3, E6, E12, E24, etc. series. The units used to specify capacitor values includes everything from picofarad (pF), nanofarad (nF), microfarad (μF) and farad (F). The percentage of allowed deviation of the capacitance from the rated value is called capacitance tolerance. The actual capacitance value must be within the tolerance limits, or the capacitor is out of specification. For abbreviated marking in tight spaces, a letter code for each tolerance is specified in IEC/EN 60062. The required capacitance tolerance is determined by the particular application. The narrow tolerances of E24 to E96 will be used for high-quality class 1 capacitors in circuits such as precision oscillators and timers. For applications such as non-critical filtering or coupling circuits, for class 2 capacitors the tolerance series E12 down to E3 are sufficient. Temperature dependence of capacitance Capacitance of ceramic capacitors varies with temperature. The different dielectrics of the many capacitor types show great differences in temperature dependence. The temperature coefficient is expressed in parts per million (ppm) per degree Celsius for class 1 ceramic capacitors or in percent (%) over the total temperature range for class 2 capacitors. Frequency dependence of capacitance Most discrete capacitor types have greater or smaller capacitance changes with increasing frequencies. The dielectric strength of class 2 ceramic and plastic film diminishes with rising frequency. Therefore, their capacitance value decreases with increasing frequency. This phenomenon is related to the dielectric relaxation in which the time constant of the electrical dipoles is the reason for the frequency dependence of permittivity. The graph on the right hand side shows typical frequency behavior for class 2 vs class 1 capacitors. Voltage dependence of capacitance Capacitance of ceramic capacitors may also change with applied voltage. This effect is more prevalent in class 2 ceramic capacitors. The ferroelectric material depends on the applied voltage. The higher the applied voltage, the lower the permittivity. Capacitance measured or applied with higher voltage can drop to values of −80% of the value measured with the standardized measuring voltage of 0.5 or 1.0 V. This behavior is a small source of nonlinearity in low-distortion filters and other analog applications. In audio applications this can be the reason for harmonic distortions. The voltage dependence of capacitance in the two diagrams above shows curves from ceramic capacitors with NME metallization. For capacitors with BME metallization the voltage dependence of capacitance increased significantly. Voltage proof For most capacitors, a physically conditioned dielectric strength or a breakdown voltage usually could be specified for each dielectric material and thickness. This is not possible with ceramic capacitors. The breakdown voltage of a ceramic dielectric layer may vary depending on the electrode material and the sintering conditions of the ceramic up to a factor of 10. A high degree of precision and control of process parameters is necessary to keep the scattering of electrical properties for today's very thin ceramic layers within specified limits. The voltage proof of ceramic capacitors is specified as rated voltage (UR). This is the maximum DC voltage that may be continuously applied to the capacitor up to the upper temperature limit. This guaranteed voltage proof is tested according to the voltages shown in the adjacent table. Furthermore, in periodic life time tests (endurance tests) the voltage proof of ceramic capacitors is tested with increased test voltage (120 to 150% of UR) to ensure safe construction. Impedance The frequency dependent AC resistance of a capacitor is called impedance and is a complex ratio of voltage to current in an AC circuit. Impedance extends the concept of Ohm's law to AC circuits, and possesses both magnitude and phase at a particular frequency, unlike resistance, which has only magnitude. Impedance is a measure of the ability of the capacitor to pass alternating currents. In this sense impedance can be used like Ohms law to calculate either the peak or the effective value of the current or the voltage. As shown in the series-equivalent circuit of a capacitor, the real-world component includes an ideal capacitor , an inductance and a resistor . To calculate the impedance the resistance and then both reactances have to be added geometrically wherein the capacitive reactance (Capacitance) is and an inductive reactance (Inductance) is . In the special case of resonance, in which both reactive resistances have the same value (), then the impedance will only be determined by . Data sheets of ceramic capacitors only specify the impedance magnitude . The typical impedance curve shows that with increasing frequency, impedance decreases, down to a minimum. The lower the impedance, the more easily alternating currents can pass through the capacitor. At the minimum point of the curve, the point of resonance, where XC has the same value as XL, the capacitor exhibits its lowest impedance value. Here only the ohmic ESR determines the impedance. With frequencies above the resonance, impedance increases again due to the ESL. ESR, dissipation factor, and quality factor The summarized losses in ceramic capacitors are ohmic AC losses. DC losses are specified as "leakage current" or "insulating resistance" and are negligible for an AC specification. These AC losses are nonlinear and may depend on frequency, temperature, age, and for some special types, on humidity. The losses result from two physical conditions, line losses with internal supply line resistances, the contact resistance of the electrode contact, the line resistance of the electrodes the dielectric losses out of the dielectric polarization The largest share of these losses in larger capacitors is usually the frequency dependent ohmic dielectric losses. Regarding the IEC 60384-1 standard, the ohmic losses of capacitors are measured at the same frequency used to measure capacitance. These are: 100 kHz, 1 MHz (preferred) or 10 MHz for ceramic capacitors with CR ≤ 1 nF: 1 kHz or 10 kHz for ceramic capacitors with 1 nF < CR ≤ 10 μF 50/60 Hz or 100/120 Hz for ceramic capacitors with CR > 10 μF Results of the summarized resistive losses of a capacitor may be specified either as equivalent series resistance (ESR), as dissipation factor (DF, tan δ), or as quality factor (Q), depending on the application requirements. Class 2 capacitors are mostly specified with the dissipation factor, tan δ. The dissipation factor is determined as the tangent of the reactance – and the ESR, and can be shown as the angle δ between the imaginary and impedance axes in the above vector diagram, see paragraph "Impedance". If the inductance is small, the dissipation factor can be approximated as: Class 1 capacitors with very low losses are specified with a dissipation factor and often with a quality factor (Q). The quality factor is defined as the reciprocal of the dissipation factor. The Q factor represents the effect of electrical resistance, and characterizes a resonator's bandwidth relative to its center or resonant frequency . A high Q value is a mark of the quality of the resonance for resonant circuits. In accordance with IEC 60384-8/-21/-9/-22 ceramic capacitors may not exceed the following dissipation factors: The ohmic losses of ceramic capacitors are frequency, temperature and voltage dependent. Additionally, class 2 capacitor measurements change because of aging. Different ceramic materials have differing losses over the temperature range and the operating frequency. The changes in class 1 capacitors are in the single-digit range while class 2 capacitors have much higher changes. HF use, inductance (ESL) and self-resonant frequency Electrical resonance occurs in a ceramic capacitor at a particular resonance frequency where the imaginary parts of the capacitor impedance and admittances cancel each other. This frequency where XC is as high as XL is called the self-resonant frequency and can be calculated with: where ω = 2πf, in which f is the resonance frequency in Hertz, L is the inductance in henries, and C is the capacitance in farads. The smaller the capacitance C and the inductance L the higher is the resonance frequency. The self-resonant frequency is the lowest frequency at which impedance passes through a minimum. For any AC application the self-resonant frequency is the highest frequency at which a capacitor can be used as a capacitive component. At frequencies above the resonance, the impedance increases again due to ESL: the capacitor becomes an inductor with inductance equal to capacitor's ESL, and resistance equal to ESR at the given frequency. ESL in industrial capacitors is mainly caused by the leads and internal connections used to connect the plates to the outside world. Larger capacitors tend to higher ESL than small ones, because the distances to the plate are longer and every millimeter increases inductance. Ceramic capacitors, which are available in the range of very small capacitance values (pF and higher) are already out of their smaller capacitance values suitable for higher frequencies up to several 100 MHz (see formula above). Due to the absence of leads and proximity to the electrodes, MLCC chips have significantly lower parasitic inductance than f. e. leaded types, which makes them suitable for higher frequency applications. A further reduction of parasitic inductance is achieved by contacting the electrodes on the longitudinal side of the chip instead of the lateral side. Sample self-resonant frequencies for one set of NP0/C0G and one set of X7R ceramic capacitors are: Note that X7Rs have better frequency response than C0Gs. It makes sense, however, since class 2 capacitors are much smaller than class 1, so they ought to have lower parasitic inductance. Aging In ferroelectric class 2 ceramic capacitors capacitance decreases over time. This behavior is called "aging". Aging occurs in ferroelectric dielectrics, where domains of polarization in the dielectric contribute to total polarization. Degradation of the polarized domains in the dielectric decreases permittivity over time so that the capacitance of class 2 ceramic capacitors decreases as the component ages. The aging follows a logarithmic law. This law defines the decrease of capacitance as a percentage for a time decade after the soldering recovery time at a defined temperature, for example, in the period from 1 to 10 hours at 20 °C. As the law is logarithmic, the percentage loss of capacitance will twice between 1 h and 100 h and 3 times between 1 h and 1000 h and so on. So aging is fastest near the beginning, and the capacitance value effectively stabilizes over time. The rate of aging of class 2 capacitors mainly depends on the materials used. A rule of thumb is, the higher the temperature dependence of the ceramic, the higher the aging percentage. The typical aging of X7R ceramic capacitors is about 2.5% per decade The aging rate of Z5U ceramic capacitors is significantly higher and can be up to 7% per decade. The aging process of class 2 capacitors may be reversed by heating the component above the Curie point. Class 1 capacitors do not experience ferroelectric aging like Class 2's. But environmental influences such as higher temperature, high humidity and mechanical stress can, over a longer period of time, lead to a small irreversible decline in capacitance, sometimes also called aging. The change of capacitance for P 100 and N 470 Class 1's is lower than 1%, for capacitors with N 750 to N 1500 ceramics it is ≤ 2%. Insulation resistance and self-discharge constant The resistance of the dielectric is never infinite, leading to some level of DC "leakage current", which contributes to self-discharge. For ceramic capacitors this resistance, placed in parallel with the capacitor in the series-equivalent circuit of capacitors, is called "insulation resistance Rins". The insulation resistance must not be confused with the outer isolation with respect to the environment. The rate of self-discharge with decreasing capacitor voltage follows the formula With the stored DC voltage and the self-discharge constant That means, after capacitor voltage dropped to 37% of the initial value. The insulation resistance given in the unit MΩ (106 Ohm) as well as the self-discharge constant in seconds is an important parameter for the quality of the dielectric insulation. These time values are important, for example, when a capacitor is used as timing component for relays or for storing a voltage value as in a sample and hold circuits or operational amplifiers. In accordance with the applicable standards, Class 1 ceramic capacitors have an Rins ≥ 10,000 MΩ for capacitors with CR ≤ 10 nF or τs ≥ 100 s for capacitors with CR > 10 nF. Class 2 ceramic capacitors have an Rins ≥ 4,000 MΩ for capacitors with CR ≤ 25 nF or τs ≥ 100 s for capacitors with CR > 25 nF. Insulation resistance and thus the self-discharge time rate are temperature dependent and decrease with increasing temperature at about 1 MΩ per 60 °C. Dielectric absorption (soakage) Dielectric absorption is the name given to the effect by which a capacitor, which has been charged for a long time, discharges only incompletely. Although an ideal capacitor remains at zero volts after discharge, real capacitors will develop a small voltage coming from time-delayed dipole discharging, a phenomenon that is also called dielectric relaxation, "soakage" or "battery action". In many applications of capacitors dielectric absorption is not a problem but in some applications, such as long-time-constant integrators, sample-and-hold circuits, switched-capacitor analog-to-digital converters and very low-distortion filters, it is important that the capacitor does not recover a residual charge after full discharge, and capacitors with low absorption are specified. The voltage at the terminals generated by dielectric absorption may in some cases possibly cause problems in the function of an electronic circuit or can be a safety risk to personnel. To prevent shocks, most very large capacitors like power capacitors are shipped with shorting wires that are removed before use. Microphony All class 2 ceramic capacitors using ferroelectric ceramics exhibit piezoelectricity, and have a piezoelectric effect called microphonics, microphony or in audio applications squealing. Microphony describes the phenomenon wherein electronic components transform mechanical vibrations into an electrical signal which in many cases is undesired noise. Sensitive electronic preamplifiers generally use class 1 ceramic and film capacitors to avoid this effect. In the reverse microphonic effect, the varying electric field between the capacitor plates exerts a physical force, moving them as a speaker. High current impulse loads or high ripple currents can generate audible acoustic sound coming from the capacitor, but discharges the capacitor and stresses the dielectric. Soldering Ceramic capacitors may experience changes to their electrical parameters due to soldering stress. The heat of the solder bath, especially for SMD styles, can cause changes of contact resistance between terminals and electrodes. For ferroelectric class 2 ceramic capacitors, the soldering temperature is above the Curie point. The polarized domains in the dielectric are going back and the aging process of class 2 ceramic capacitors is starting again. Hence after soldering a recovery time of approximately 24 hours is necessary. After recovery some electrical parameters like capacitance value, ESR, leakage currents are changed irreversibly. The changes are in the lower percentage range depending on the style of capacitor. Additional information Standardization The standardization for all electrical, electronic components and related technologies follows the rules given by the International Electrotechnical Commission (IEC), a non-profit, non-governmental international standards organization. The definition of the characteristics and the procedure of the test methods for capacitors for use in electronic equipment are set out in the generic specification: IEC 60384-1, Fixed capacitors for use in electronic equipment – Part 1: Generic specification The tests and requirements to be met by ceramic capacitors for use in electronic equipment for approval as standardized types are set out in the following sectional specifications: IEC 60384-8, Fixed capacitors of ceramic dielectric, Class 1 IEC 60384-9, Fixed capacitors of ceramic dielectric, Class 2 IEC 60384-21, Fixed surface mount multilayer capacitors of ceramic dielectric, Class 1 IEC 60384-22, Fixed surface mount multilayer capacitors of ceramic dielectric, Class 2 Tantalum capacitor replacement Multilayer ceramic capacitors are increasingly used to replace tantalum and low capacitance aluminium electrolytic capacitors in applications such as bypass or high frequency switched-mode power supplies as their cost, reliability and size becomes competitive. In many applications, their low ESR allows the use of a lower nominal capacitance value. Features and disadvantages of ceramic capacitors For features and disadvantages of ceramic capacitors see main article Capacitor types#Comparison of types Marking Imprinted markings If space permits ceramic capacitors, like most other electronic components, have imprinted markings to indicate the manufacturer, the type, their electrical and thermal characteristics and their date of manufacture. In the ideal case, if they are large enough, the capacitor will be marked with: manufacturer's name or trademark; manufacturer's type designation; rated capacitance; tolerance on rated capacitance rated voltage and nature of supply (AC or DC) climatic category or rated temperature; year and month (or week) of manufacture; certification marks of safety standards (for safety EMI/RFI suppression capacitors) Smaller capacitors use a shorthand notation, to display all the relevant information in the limited space. The most commonly used format is: XYZ J/K/M VOLTS V, where XYZ represents the capacitance (calculated as XY × 10Z pF), the letters J, K or M indicate the tolerance (±5%, ±10% and ±20% respectively) and VOLTS V represents the working voltage. Examples A capacitor with the following text on its body: 105K 330V has a capacitance of 10 × 105 pF = 1 μF (K = ±10%) with a working voltage of 330 V. A capacitor with the following text: 473M 100V has a capacitance of 47 × 103 pF = 47 nF (M = ±20%) with a working voltage of 100 V. Capacitance, tolerance and date of manufacture can be identified with a short code according to IEC/EN 60062. Examples of short-marking of the rated capacitance (microfarads): μ47 = 0.47 μF 4μ7 = 4.7 μF 47μ = 47 μF The date of manufacture is often printed in accordance with international standards. Version 1: coding with year/week numeral code, "1208" is "2012, week number 8". Version 2: coding with year code/month code, Year code: "R" = 2003, "S"= 2004, "T" = 2005, "U" = 2006, "V" = 2007, "W" = 2008, "X" = 2009, "A" = 2010, "B" = 2011, "C" = 2012, "D" = 2013 etc. Month code: "1" to "9" = Jan. to Sept., "O" = October, "N" = November, "D" = December "X5" is then "2009, May" For very small capacitors like MLCC chips no marking is possible. Here only the traceability of the manufacturers can ensure the identification of a type. Colour coding The identification of modern capacitors has no detailed color coding.	Technology	Components	null
6935153	https://en.wikipedia.org/wiki/Cuvier%27s%20dwarf%20caiman	Cuvier's dwarf caiman	Cuvier's dwarf caiman (Paleosuchus palpebrosus) is a small crocodilian in the alligator family from northern and central South America. It is found in Bolivia, Brazil, Colombia, Ecuador, French Guiana, Guyana, Paraguay, Peru, Suriname, Trinidad and Venezuela. It lives in riverine forests, flooded forests near lakes, and near fast-flowing rivers and streams. It can traverse dry land to reach temporary pools and tolerates colder water than other species of caimans. Other common names for this species include the musky caiman, the dwarf caiman, Cuvier's caiman, and the smooth-fronted caiman (the latter name is also used for P. trigonatus). It is sometimes kept in captivity as a pet and may be referred to as the wedge-head caiman by the pet trade community. Cuvier's dwarf caiman was first described by the French zoologist Georges Cuvier in 1807 and is one of only two species in the genus Paleosuchus, the other species being P. trigonatus. Their closest relatives are the other caimans in the subfamily Caimaninae. With a total length averaging for males and up to for females, Cuvier's dwarf caiman is not only the smallest extant species in the alligator and caiman family, but also the smallest of all crocodilians (unless the Congo dwarf crocodile is considered a valid species). An adult weighs around . Its lack of size is partly made up for by its strong body armor, provided by the bony bases to its dermal scales, which provides protection against predators. Juvenile dwarf caimans mainly feed on invertebrates, but also small fish and frogs, while adults eat larger fish, amphibians, and invertebrates, such as large molluscs. This caiman sometimes uses a burrow as shelter during the day and in the Pantanal may aestivate in the burrow to stay cool in the dry season. The female buries her eggs on a mounded nest and these take about 3 months to hatch. She helps the hatchlings to escape from the nest and provides some parental care for the first few weeks of their lives. This caiman has a wide range and large total population and the IUCN lists its conservation status as being of least concern. Etymology The genus name Paleosuchus is derived from the Greek palaios meaning "ancient" and soukhos meaning "crocodile". This refers to the belief that this crocodile comes from an ancient lineage that diverged from other species of caimans some 30 million years ago. The specific name palpebrosus is derived from the Latin palpebra meaning "eyelid" and osus meaning "full of". This refers to the bony plates (palpebrals) present on the upper eyelids. Common names include the musky caiman, the dwarf caiman, Cuvier's caiman, and the smooth-fronted caiman, although the last of these is also used to refer to the closely related P. trigonatus. In the pet trade, it is sometimes referred to as the wedge-head caiman. Discovery and taxonomy Cuvier's dwarf caiman was first described by Cuvier in 1807 as Crocodylus palpebrosus from a type locality described as "Cayenne". Since then, it has been given a number of names by different authorities: Crocodilus (Alligator) palpebrosus (Merrem, 1820), Jacaretinga moschifer (Spix, 1825), Champsa palpebrosa (Wagler, 1830), Alligator palpebrosus (Dumeril and Bibron, 1836), Champsa gibbiceps (Natterer, 1841), Caiman palpebrosus (Gray, 1844), Caiman (Aromosuchus) palpebrosus (Gray, 1862), and Jacaretinga palpebrosus (Vaillant, 1898). Muller, in 1924, and Schmidt, in 1928, were the first to use the currently accepted name of Paleosuchus palpebrosus. No subspecies are recognised. At present, the genus Paleosuchus contains only two members, Paleosuchus trigonatus, commonly known as the smooth-fronted or Schneider's dwarf caiman, and P. palpebrosus, both from South America. Paleosuchus is distinguished from other caimans in the alligator subfamily Caimaninae by the absence of an interorbital ridge and the presence of four teeth in the premaxilla region of the jaw, where other species of caimans have five. The relationships of extant (living) caimans can be shown in the cladogram below, based on molecular DNA-based phylogenetic studies: A genetic study in 2012 found clear differences between various populations of Cuvier's dwarf caiman (Pantanal; Madeira River basin; Rio Negro basin), and these are apparently isolated from each other, leading to the suggestion that it may be a cryptic species complex. Description Cuvier's dwarf caiman is the smallest living New World crocodilian. Males grow to a maximum length around while females do not usually exceed . The largest specimen on record measured in length. This may be an underestimate of the animal's maximum size, as nearly all large adults have lost the tips of their tails and the largest specimen measured in the Pantanal region had a snout–vent length of (equivalent to a total length of with an intact tail). An adult typically weighs around , around the same weight as a 6- to 12-month-old specimen of several larger species of crocodilians. Large adults of this species can weigh up to . Cuvier's dwarf caiman has strong body armor on both its dorsal (upper) and ventral (lower) sides, which may compensate for its small body size in reducing predation. The dermal scales providing this protection have a bony base and are known as osteoderms. The head has an unusual shape for a crocodilian, with a dome-shaped skull and a short, smooth, concave snout with an upturned tip, the shape rather resembling the head of a dog. The upper jaw extends markedly further forward than the lower jaw. Four premaxillary and 14 to 15 maxillary teeth are on either side of the upper jaw and 21 or 22 teeth on each side of the lower jaw, giving a total of about 80 teeth. The neck is relatively slender and the dorsal scutes are less prominent than in the smooth-fronted caiman. The double rows of scutes on the tail are small and project vertically. Adults are dark brownish-black with a dark brown head, while juveniles are brown with black bands. The irises of the eyes are chestnut brown at all ages and the pupils are vertical slits.The scutellation (arrangement of the scales) helps to distinguish Cuvier's dwarf caiman from Schneider's dwarf caiman. Distribution and habitat Cuvier's dwarf caiman is native to tropical northern and central South America. It is present in the drainages of the Orinoco River, the São Francisco River, and the Amazon River, and the upper reaches of the Paraná River and the Paraguay River. The countries in which it is found include Peru, Ecuador, Colombia, Venezuela, Guyana, Suriname, French Guiana, Brazil, Bolivia, Trinidad and Paraguay. The range of this species is rather larger than that of the sympatric smooth-fronted caiman, as it extends into Paraguay and includes a larger area of Brazil. They also follow seasonal fluctuations in water-level, while the smooth-fronted caiman does not, which may explain how the two species are able to live in sympatry. Cuvier's dwarf caiman is a freshwater species and is found in forested riverine habitats and areas of flooded forest around lakes. It seems to prefer rivers and streams with fast-flowing water, but it is also found in quiet, nutrient-poor waters in Venezuela and southeastern Brazil. It is able to travel quite large distances overland at night and subadult individuals have sometimes been found in isolated, temporary pools. In the northern and southern parts of its range, it is also found in gallery forests in savanna country, but it is absent from such habitats in the Llanos and the Pantanal. Cuvier's dwarf caiman seems relatively tolerant of cool water compared to other species of caimans. During the day, individuals sometimes lie up in burrows but at other times rest on piles of rocks or sun themselves while lying, facing the sun, in shallow water with their backs exposed. Behaviour and ecology These caimans are mainly nocturnal. Because they occupy many different microhabitats, their diet is believed to vary regionally. Adults feed on fish, amphibians, small mammals, birds, crabs, shrimp, molluscs, insects, and other invertebrates, which they catch in the water or on land. Juveniles eat fewer fish, but also consume crustaceans, tadpoles, frogs, and snails, as well as land invertebrates, such as beetles. The prey is mostly swallowed whole and is ground up by stones in the gizzard. In the Pantanal, Cuvier's dwarf caiman estivates in burrows during the dry season and is able to maintain its temperature around for days at a time. Adult Cuvier's dwarf caimans are usually found singly or in pairs. The breeding of this species has been little studied, but it does not appear to be seasonal in nature. The female builds a mound nest out of vegetation and mud somewhere in a concealed location and lays a clutch of 10 to 25 eggs, hiding them under further vegetation. Nest temperature varies between and are heated by decaying vegetation. The incubation period is around 90 days and the sex of the hatchlings depends on the temperature of the nest during that time. When the eggs begin to hatch, the female opens the nest in response to the calls made by the young. Newly emerged juveniles have a coating of mucus and may delay entering the water for a few days until this has dried. Its continuing presence on their skin is believed to reduce algal growth. The female stays with the young for around a year, with the longest recorded care extending to 21 months. After this the hatchlings disperse. The young grow at a rate around per year. Females reach sexual maturity around 8 years old and males around 6 years old. Cuvier's dwarf caiman is considered to be a keystone species whose presence in the ecosystem maintains a healthy balance of organisms. In its absence, fish, such as piranhas, might dominate the environment. The eggs and newly hatched young are most at risk and are preyed on by birds, snakes, rats, raccoons, and other mammals. Adults are protected by the bony osteoderms under the scales and their main predators are jaguars, green anacondas (Eunectes murinus), and large boa constrictors (Boa constrictor). The Cuvier's dwarf caiman is the only crocodilian species that seemingly does not perform the near-universal "death roll" technique used by other extant crocodilians for feeding or intra-specific combat. However, this may only be circumstantial, as specimens tested for the behavior may have been acting uncooperatively with the researchers. Status and conservation Many crocodilians are hunted for their skins, but this is not the case with the Cuvier's dwarf caiman. This may be because the ventral skin in this species is too heavily armored to make it easy to tan. Some individuals are killed by indigenous peoples for food and some traditional South American cultures believe dwarf caiman teeth protect from snake bites. Others, particularly in Guyana, are collected for the pet trade; but no evidence shows that populations are dwindling as a result. Some threats to this species are from habitat destruction, including the mining of gold, but these are not thought to be of great significance. The estimated total population is over a million individuals. In its Red List of Threatened Species, the IUCN lists Cuvier's dwarf caiman as being of least concern, which is because its range is extensive, covering much of northern and central South America, and although its population trend is unknown, it appears to be abundant in many of the localities in which it is found. It is listed in Appendix II of CITES. Captive care Cuvier's dwarf caiman can be kept as a pet, though providing suitable care is expensive and requires extremely large enclosures. In many countries, permits or licenses are necessary and most veterinarians have little experience with these exotic animals.	Biology and health sciences	Crocodilia	Animals
10725984	https://en.wikipedia.org/wiki/Gunshot%20wound	Gunshot wound	A gunshot wound (GSW) is a penetrating injury caused by a projectile (e.g. a bullet) shot from a gun (typically a firearm). Damage may include bleeding, bone fractures, organ damage, wound infection, and loss of the ability to move part of the body. Damage depends on the part of the body hit, the path the bullet follows through (or into) the body, and the type and speed of the bullet. In severe cases, although not uncommon, the injury is fatal. Long-term complications can include bowel obstruction, failure to thrive, neurogenic bladder and paralysis, recurrent cardiorespiratory distress and pneumothorax, hypoxic brain injury leading to early dementia, amputations, chronic pain and pain with light touch (hyperalgesia), deep venous thrombosis with pulmonary embolus, limb swelling and debility, and lead poisoning. Factors that determine rates of gun violence vary by country. These factors may include the illegal drug trade, easy access to firearms, substance misuse including alcohol, mental health problems, firearm laws, social attitudes, economic differences, and occupations such as being a police officer. Where guns are more common, altercations more often end in death. Before management begins, the area must be verified as safe. This is followed by stopping major bleeding, then assessing and supporting the airway, breathing, and circulation. Firearm laws, particularly background checks and permit to purchase, decrease the risk of death from firearms. Safer firearm storage may decrease the risk of firearm-related deaths in children. In 2015, about a million gunshot wounds occurred from interpersonal violence. In 2016, firearms resulted in 251,000 deaths globally, up from 209,000 in 1990. Of these deaths, 161,000 (64%) were the result of assault, 67,500 (27%) were the result of suicide, and 23,000 (9%) were accidents. In the United States, guns resulted in about 40,000 deaths in 2017. Firearm-related deaths are most common in males between the ages of 20 and 24 years. Economic costs due to gunshot wounds have been estimated at $140 billion a year in the United States. Signs and symptoms Trauma from a gunshot wound varies widely based on the bullet, velocity, mass, entry point, trajectory, affected anatomy, and exit point. Gunshot wounds can be particularly devastating compared to other penetrating injuries because the trajectory and fragmentation of bullets can be unpredictable after entry. Moreover, gunshot wounds typically involve a large degree of nearby tissue disruption and destruction caused by the physical effects of the projectile correlated with the bullet velocity classification. The immediate damaging effect of a gunshot wound is typically severe bleeding with the potential for a type of shock known as hypovolemic shock, a condition characterized by inadequate delivery of oxygen to vital organs. In the case of traumatic hypovolemic shock, this failure of adequate oxygen delivery is due to blood loss, as blood is the means of delivering oxygen to the body's constituent parts. Besides blood loss, internal bleeding can lead to complications. Devastating effects can result when a bullet strikes a vital organ such as the heart, lungs, or liver, or damages a component of the central nervous system such as the spinal cord or brain. It can lead to organ failure and death. Common causes of death following gunshot injury include bleeding, low oxygen caused by pneumothorax, catastrophic injury to the heart and major blood vessels, and damage to the brain or central nervous system. Non-fatal gunshot wounds frequently have mild to severe long-lasting effects, typically some form of major disfigurement such as amputation because of a severe bone fracture and may cause permanent disability. A sudden blood gush may take effect immediately from a gunshot wound if a bullet directly damages larger blood vessels, especially arteries. Pathophysiology The degree of tissue disruption caused by a projectile is related to the cavitation the projectile creates as it passes through tissue. A bullet with sufficient energy will have a cavitation effect in addition to the penetrating track injury. As the bullet passes through the tissue, initially crushing then lacerating, the space left forms a cavity; this is called the permanent cavity. Higher-velocity bullets create a pressure wave that forces the tissues away, creating not only a permanent cavity the size of the caliber of the bullet but a temporary cavity or secondary cavity, which is often many times larger than the bullet itself. The temporary cavity is the radial stretching of tissue around the bullet's wound track, which momentarily leaves an empty space caused by high pressures surrounding the projectile that accelerate material away from its path. The extent of cavitation, in turn, is related to the following characteristics of the projectile: Kinetic energy: KE = 1/2mv2 (where m is mass and v is velocity). This helps to explain why wounds produced by projectiles of higher mass and/or higher velocity produce greater tissue disruption than projectiles of lower mass and velocity. The velocity of the bullet is a more important determinant of tissue injury. Although both mass and velocity contribute to the overall energy of the projectile, the energy is proportional to the mass while proportional to the square of its velocity. As a result, for constant velocity, if the mass is doubled, the energy is doubled; however, if the velocity of the bullet is doubled, the energy increases four times. The initial velocity of a bullet is largely dependent on the firearm. The US military commonly uses 5.56-mm bullets, which have a relatively low mass as compared with other bullets; however, the speed of these bullets is relatively fast. As a result, they produce a larger amount of kinetic energy, which is transmitted to the tissues of the target. The size of the temporary cavity is approximately proportional to the kinetic energy of the bullet and depends on the resistance of the tissue to stress. Muzzle energy, which is based on muzzle velocity, is often used for ease of comparison. Yaw: Handgun bullets will generally travel in a relatively straight line or make one turn if a bone is hit. Upon travel through deeper tissue, high-energy rounds may become unstable as they decelerate, and may tumble (pitch and yaw) as the energy of the projectile is absorbed, causing stretching and tearing of the surrounding tissue. Fragmentation: Most commonly, bullets do not fragment, and secondary damage from fragments of shattered bone is a more common complication than bullet fragments. Diagnosis Classification Gunshot wounds are classified according to the speed of the projectile using the Gustilo open fracture classification: Low-velocity: Less than 335 m/s (1,100 ft/s) Low velocity wounds are typical of small caliber handguns. They do not usually cause extensive soft tissue damage, and in the Gustilo open fracture classification are classified as Type 1 or 2 wounds. Medium-velocity: Between 360 m/s (1,200 ft/s) and 600 m/s (2,000 ft/s) These are more typical of shotgun blasts or higher caliber handguns like magnums. The risk of infection from these types of wounds can vary depending on the type and pattern of bullets fired as well as the distance from the firearm. High-velocity: Between 600 m/s (2,000 ft/s) and 1,000 m/s (3,500 ft/s) Usually caused by powerful assault or hunting rifles and usually cause Gustilo Type 3 wounds. The risk of infection is especially high due to the large area of injury and destroyed tissue. Bullets from handguns are sometimes less than but with modern pistol loads, they usually are slightly above , while bullets from most modern rifles exceed . One recently developed class of firearm projectiles is the hyper-velocity bullet, such cartridges are usually made for achieving such high speed, purpose-built in factories or made by amateurs. Examples of hyper velocity cartridges include the .220 Swift, .17 Remington and .17 Mach IV cartridges. The US military commonly uses 5.56mm bullets, which have a relatively low mass as compared with other bullets (2,6-4,0 grams); however, the speed of these bullets is relatively fast (approximately , placing them in the high velocity category). As a result, they produce a larger amount of kinetic energy, which is transmitted to the tissues of the target. High energy transfer results in more tissue disruption, which plays a role in incapacitation, but other factors such as wound size and shot placement are also important. Kronlein shot The "Kronlein shot" (German: Krönleinschuss) is a distinctive type of headshot wound that can only be created by a high velocity rifle bullet or shotgun slug. In a Kronlein shot, the intact brain is ejected from the skull and deposited some distance from the victim's body. This type of wound is believed to be caused by a hydrodynamic effect. Hydraulic pressure generated within the skull by a high velocity bullet leads to the explosive ejection of the brain from the fractured skull. Prevention Interventions have been recommended to reduce the risk of firearm related injury or death. Medical organizations in the United States recommend a criminal background check being held before a person buys a gun and that a person who has convictions for crimes of violence should not be permitted to buy a gun. Safe storage of guns is recommended, as well as better mental health care and removal of guns from those at risk of suicide. Experts recommend that physicians counsel patients regarding safe storage of guns and other injury prevention strategies related to guns as part of routine medical care. Having guns locked and unloaded is associated with a lower risk of gun related injury or death (including a lower risk of suicide) for all household members as compared to guns that are stored loaded and unlocked. Temporarily removing guns from the home, either voluntarily or by court order (such as with extreme risk protection orders [so called "red flag laws"] in the United States) is recommended for those who are at risk of suicide or violence towards others. Such laws have been associated with a lower risk of suicide using guns in population based studies. In an effort to prevent mass shootings, greater regulations on guns that can rapidly fire many bullets is recommended. Management Initial assessment for a gunshot wound is approached in the same way as other acute trauma using the advanced trauma life support (ATLS) protocol. These include: A) Airway - Assess and protect airway and potentially the cervical spine B) Breathing - Maintain adequate ventilation and oxygenation C) Circulation - Assess for and control bleeding to maintain organ perfusion including focused assessment with sonography for trauma (FAST) D) Disability - Perform basic neurological exam including Glasgow Coma Scale (GCS) E) Exposure - Expose entire body and search for any missed injuries, entry points, and exit points while maintaining body temperature Depending on the extent of injury, management can range from urgent surgical intervention to observation. As such, any history from the scene such as gun type, shots fired, shot direction and distance, blood loss on scene, and pre-hospital vitals signs can be very helpful in directing management. Unstable people with signs of bleeding that cannot be controlled during the initial evaluation require immediate surgical exploration in the operating room. Otherwise, management protocols are generally dictated by anatomic entry point and anticipated trajectory. Neck A gunshot wound to the neck can be particularly dangerous because of the high number of vital anatomical structures contained within a small space. The neck contains the larynx, trachea, pharynx, esophagus, vasculature (carotid, subclavian, and vertebral arteries; jugular, brachiocephalic, and vertebral veins; thyroid vessels), and nervous system anatomy (spinal cord, cranial nerves, peripheral nerves, sympathetic chain, brachial plexus). Gunshots to the neck can thus cause severe bleeding, airway compromise, and nervous system injury. Initial assessment of a gunshot wound to the neck involves non-probing inspection of whether the injury is a penetrating neck injury (PNI), classified by violation of the platysma muscle. If the platysma is intact, the wound is considered superficial and only requires local wound care. If the injury is a PNI, surgery should be consulted immediately while the case is being managed. Of note, wounds should not be explored on the field or in the emergency department given the risk of exacerbating the wound. Due to the advances in diagnostic imaging, management of PNI has been shifting from a "zone-based" approach, which uses anatomical site of injury to guide decisions, to a "no-zone" approach which uses a symptom-based algorithm. The no-zone approach uses a hard signs and imaging system to guide next steps. Hard signs include airway compromise, unresponsive shock, diminished pulses, uncontrolled bleeding, expanding hematoma, bruits/thrill, air bubbling from wound or extensive subcutaneous air, stridor/hoarseness, neurological deficits. If any hard signs are present, immediate surgical exploration and repair is pursued alongside airway and bleeding control. If there are no hard signs, the person receives a multi-detector CT angiography for better diagnosis. A directed angiography or endoscopy may be warranted in a high-risk trajectory for the gunshot. A positive finding on CT leads to operative exploration. If negative, the person may be observed with local wound care. Chest Important anatomy in the chest includes the chest wall, ribs, spine, spinal cord, intercostal neurovascular bundles, lungs, bronchi, heart, aorta, major vessels, esophagus, thoracic duct, and diaphragm. Gunshots to the chest can thus cause severe bleeding (hemothorax), respiratory compromise (pneumothorax, hemothorax, pulmonary contusion, tracheobronchial injury), cardiac injury (pericardial tamponade), esophageal injury, and nervous system injury. Initial workup as outlined in the Workup section is particularly important with gunshot wounds to the chest because of the high risk for direct injury to the lungs, heart, and major vessels. Important notes for the initial workup specific for chest injuries are as follows. In people with pericardial tamponade or tension pneumothorax, the chest should be evacuated or decompressed if possible prior to attempting tracheal intubation because the positive pressure ventilation can cause hypotention or cardiovascular collapse. Those with signs of a tension pneumothorax (asymmetric breathing, unstable blood flow, respiratory distress) should immediately receive a chest tube (> French 36) or needle decompression if chest tube placement is delayed. FAST exam should include extended views into the chest to evaluate for hemopericardium, pneumothorax, hemothorax, and peritoneal fluid. Those with cardiac tamponade, uncontrolled bleeding, or a persistent air leak from a chest tube all require surgery. Cardiac tamponade can be identified on FAST exam. Blood loss warranting surgery is 1–1.5 L of immediate chest tube drainage or ongoing bleeding of 200-300 mL/hr. Persistent air leak is suggestive of tracheobronchial injury which will not heal without surgical intervention. Depending on the severity of the person's condition and if cardiac arrest is recent or imminent, the person may require surgical intervention in the emergency department, otherwise known as an emergency department thoracotomy (EDT). However, not all gunshot to the chest require surgery. Asymptomatic people with a normal chest X-ray can be observed with a repeat exam and imaging after 6 hours to ensure no delayed development of pneumothorax or hemothorax. If a person only has a pneumothorax or hemothorax, a chest tube is usually sufficient for management unless there is large volume bleeding or persistent air leak as noted above. Additional imaging after initial chest X-ray and ultrasound can be useful in guiding next steps for stable people. Common imaging modalities include chest CT, formal echocardiography, angiography, esophagoscopy, esophagography, and bronchoscopy depending on the signs and symptoms. Abdomen Important anatomy in the abdomen includes the stomach, small bowel, colon, liver, spleen, pancreas, kidneys, spine, diaphragm, descending aorta, and other abdominal vessels and nerves. Gunshots to the abdomen can thus cause severe bleeding, release of bowel contents, peritonitis, organ rupture, respiratory compromise, and neurological deficits. The most important initial evaluation of a gunshot wound to the abdomen is whether there is uncontrolled bleeding, inflammation of the peritoneum, or spillage of bowel contents. If any of these are present, the person should be transferred immediately to the operating room for laparotomy. If it is difficult to evaluate for those indications because the person is unresponsive or incomprehensible, it is up to the surgeon's discretion whether to pursue laparotomy, exploratory laparoscopy, or alternative investigative tools. Although all people with abdominal gunshot wounds were taken to the operating room in the past, practice has shifted in recent years with the advances in imaging to non-operative approaches in more stable people. If the person's vital signs are stable without indication for immediate surgery, imaging is done to determine the extent of injury. Ultrasound (FAST) and help identify intra-abdominal bleeding and X-rays can help determine bullet trajectory and fragmentation. However, the best and preferred mode of imaging is high-resolution multi-detector CT (MDCT) with IV, oral, and sometimes rectal contrast. Severity of injury found on imaging will determine whether the surgeon takes an operative or close observational approach. Diagnostic peritoneal lavage (DPL) has become largely obsolete with the advances in MDCT, with use limited to centers without access to CT to guide requirement for urgent transfer for operation. Extremities The four main components of extremities are bones, vessels, nerves, and soft tissues. Gunshot wounds can thus cause severe bleeding, fractures, nerve deficits, and soft tissue damage. The Mangled Extremity Severity Score (MESS) is used to classify the severity of injury and evaluates for severity of skeletal and/or soft tissue injury, limb ischemia, shock, and age. Depending on the extent of injury, management can range from superficial wound care to limb amputation. Vital sign stability and vascular assessment are the most important determinants of management in extremity injuries. As with other traumatic cases, those with uncontrolled bleeding require immediate surgical intervention. If surgical intervention is not readily available and direct pressure is insufficient to control bleeding, tourniquets or direct clamping of visible vessels may be used temporarily to slow active bleeding. People with hard signs of vascular injury also require immediate surgical intervention. Hard signs include active bleeding, expanding or pulsatile hematoma, bruit/thrill, absent distal pulses and signs of extremity ischemia. For stable people without hard signs of vascular injury, an injured extremity index (IEI) should be calculated by comparing the blood pressure in the injured limb compared to an uninjured limb in order to further evaluate for potential vascular injury. If the IEI or clinical signs are suggestive of vascular injury, the person may undergo surgery or receive further imaging including CT angiography or conventional arteriography. In addition to vascular management, people must be evaluated for bone, soft tissue, and nerve injury. Plain films can be used for fractures alongside CTs for soft tissue assessment. Fractures must be debrided and stabilized, nerves repaired when possible, and soft tissue debrided and covered. This process can often require multiple procedures over time depending on the severity of injury. Epidemiology In 2015, about a million gunshot wounds occurred from interpersonal violence. Firearms, globally in 2016, resulted in 251,000 deaths up from 209,000 in 1990. Of these deaths 161,000 (64%) were the result of assault, 67,500 (27%) were the result of suicide, and 23,000 were accidents. Firearm related deaths are most common in males between the ages of 20 and 24 years. The countries with the greatest number of deaths from firearms are Brazil, United States, Mexico, Colombia, Venezuela, Guatemala, Bahamas and South Africa which make up just over half the total. In the United States in 2015, about half of the 44,000 people who died by suicide did so with a gun. As of 2016, the countries with the highest rates of gun violence per capita were El Salvador, Venezuela, and Guatemala with 40.3, 34.8, and 26.8 violent gun deaths per 100,000 people respectively. The countries with the lowest rates of were Singapore, Japan, and South Korea with 0.03, 0.04, and 0.05 violent gun deaths per 100,000 people respectively. Canada In 2016, about 893 people died due to gunshot wounds in Canada (2.1 per 100,000). About 80% were suicides, 12% were assaults, and 4% were accidents. United States In 2017, there were 39,773 deaths in the United States as a result gunshot wounds. Of these 60% were suicides, 37% were homicides, 1.4% were by law enforcement, 1.2% were accidents, and 0.9% were from an unknown cause. This is up from 37,200 deaths in 2016 due to a gunshot wound (10.6 per 100,000). With respect to those that pertain to interpersonal violence, it had the 31st highest rate in the world with 3.85 deaths per 100,000 people in 2016. The majority of all homicides and suicides are firearm-related, and the majority of firearm-related deaths are the result of murder and suicide. When sorted by GDP, however, the United States has a much higher violent gun death rate compared to other developed countries, with over 10 times the number of firearms assault deaths than the next four highest GDP countries combined. Gunshot violence is the third most costly cause of injury and the fourth most expensive form of hospitalization in the United States. History Until the 1880s, the standard practice for treating a gunshot wound called for physicians to insert their unsterilized fingers into the wound to probe and locate the path of the bullet. Standard surgical theory such as opening abdominal cavities to repair gunshot wounds, germ theory, and Joseph Lister's technique for antiseptic surgery using diluted carbolic acid, had not yet been accepted as standard practice. For example, sixteen doctors attended to President James A. Garfield after he was shot in 1881, and most probed the wound with their fingers or dirty instruments. Historians agree that massive infection was a significant factor in Garfield's death. At almost the same time, in Tombstone, Arizona Territory, on 13 July 1881, George E. Goodfellow performed the first laparotomy to treat an abdominal gunshot wound. Goodfellow pioneered the use of sterile techniques in treating gunshot wounds, washing the person's wound and his hands with lye soap or whisky, and his patient, unlike the President, recovered. He became America's leading authority on gunshot wounds and is credited as the United States' first civilian trauma surgeon. Mid-nineteenth-century handguns such as the Colt revolvers used during the American Civil War had muzzle velocities of just 230– /s and their powder and ball predecessors had velocities of 167 m/s or less. Unlike today's high-velocity bullets, nineteenth-century balls produced almost little or no cavitation and, being slower moving, they were liable to lodge in unusual locations at odds with their trajectory. Wilhelm Röntgen's discovery of X-rays in 1895 led to the use of radiographs to locate bullets in wounded soldiers. Survival rates for gunshot wounds improved among US military personnel during the Korean and Vietnam Wars, due in part to helicopter evacuation, along with improvements in resuscitation and battlefield medicine. Similar improvements were seen in US trauma practices during the Iraq War. Military health care providers who return to civilian practice sometimes disseminate military trauma care practices. One such practice is to transfer major trauma cases to an operating theater as soon as possible, to stop internal bleeding. Within the United States, the survival rate for gunshot wounds has increased, leading to declines in the gun death rate in states that have stable rates of gunshot hospitalizations.	Biology and health sciences	Types	Health
10727548	https://en.wikipedia.org/wiki/Evolution%20of%20mammals	Evolution of mammals	The evolution of mammals has passed through many stages since the first appearance of their synapsid ancestors in the Pennsylvanian sub-period of the late Carboniferous period. By the mid-Triassic, there were many synapsid species that looked like mammals. The lineage leading to today's mammals split up in the Jurassic; synapsids from this period include Dryolestes, more closely related to extant placentals and marsupials than to monotremes, as well as Ambondro, more closely related to monotremes. Later on, the eutherian and metatherian lineages separated; the metatherians are the animals more closely related to the marsupials, while the eutherians are those more closely related to the placentals. Since Juramaia, the earliest known eutherian, lived 160 million years ago in the Jurassic, this divergence must have occurred in the same period. After the Cretaceous–Paleogene extinction event wiped out the non-avian dinosaurs (birds being the only surviving dinosaurs) and several mammalian groups, placental and marsupial mammals diversified into many new forms and ecological niches throughout the Paleogene and Neogene, by the end of which all modern orders had appeared. The synapsid lineage became distinct from the sauropsid lineage in the late Carboniferous period, between 320 and 315 million years ago. The only living synapsids are mammals, while the sauropsids gave rise to the dinosaurs, and today's reptiles and birds along with all the extinct amniotes more closely related to them than to mammals. Primitive synapsids were traditionally called "mammal-like reptiles" or "pelycosaurs", but both are now seen as outdated and disfavored paraphyletic terms, since they were not reptiles, nor part of reptile lineage. The modern term for these is "stem mammals", and sometimes "protomammals" or "paramammals". Throughout the Permian period, the synapsids included the dominant carnivores and several important herbivores. In the subsequent Triassic period, however, a previously obscure group of sauropsids, the archosaurs, became the dominant vertebrates. The mammaliaforms appeared during this period; their superior sense of smell, backed up by a large brain, facilitated entry into nocturnal niches with less exposure to archosaur predation. (Conversely, mammaliaforms' success in these niches may have prevented archosaurs from becoming smaller or nocturnal themselves.) The nocturnal lifestyle may have contributed greatly to the development of mammalian traits such as endothermy and hair. Later in the Mesozoic, after theropod dinosaurs replaced rauisuchians as the dominant carnivores, mammals spread into other ecological niches. For example, some became aquatic, some were gliders, and some even fed on juvenile dinosaurs. Most of the evidence consists of fossils. For many years, fossils of Mesozoic mammals and their immediate ancestors were very rare and fragmentary; but, since the mid-1990s, there have been many important new finds, especially in China. The relatively new techniques of molecular phylogenetics have also shed light on some aspects of mammalian evolution by estimating the timing of important divergence points for modern species. When used carefully, these techniques often, but not always, agree with the fossil record. Although mammary glands are a signature feature of modern mammals, little is known about the evolution of lactation as these soft tissues are not often preserved in the fossil record. Most research concerning the evolution of mammals centers on the shapes of the teeth, the hardest parts of the tetrapod body. Other important research characteristics include the evolution of the middle ear bones, erect limb posture, a bony secondary palate, fur, hair, and warm-bloodedness. Definition of "mammal" While living mammal species can be identified by the presence of milk-producing mammary glands in the females, other features are required when classifying fossils, because mammary glands and other soft-tissue features are not visible in fossils. One such feature available for paleontology, shared by all living mammals (including monotremes), but not present in any of the early Triassic therapsids, is shown in Figure 1 (on the right), namely: mammals use two bones for hearing that all other amniotes use for eating. The earliest amniotes had a jaw joint composed of the articular (a small bone at the back of the lower jaw) and the quadrate (a small bone at the back of the upper jaw). All non-mammalian tetrapods use this system including amphibians, turtles, lizards, snakes, crocodilians, dinosaurs (including the birds), ichthyosaurs, pterosaurs and therapsids. But mammals have a different jaw joint, composed only of the dentary (the lower jaw bone, which carries the teeth) and the squamosal (another small skull bone). In the Jurassic, their quadrate and articular bones evolved into the incus and malleus bones in the middle ear. Mammals also have a double occipital condyle; they have two knobs at the base of the skull that fit into the topmost neck vertebra, while other tetrapods have a single occipital condyle. In a 1981 article, Kenneth A. Kermack and his co-authors argued for drawing the line between mammals and earlier synapsids at the point where the mammalian pattern of molar occlusion was being acquired and the dentary-squamosal joint had appeared. The criterion chosen, they noted, is merely a matter of convenience; their choice was based on the fact that "the lower jaw is the most likely skeletal element of a Mesozoic mammal to be preserved." Today, most paleontologists consider that animals are mammals if they satisfy this criterion. The ancestry of mammals Amniotes The first fully terrestrial vertebrates were reptilian amniotes — their eggs had internal membranes that allowed the developing embryo to breathe but kept water in. This allowed amniotes to lay eggs on dry land, while amphibians generally need to lay their eggs in water (a few amphibians, such as the common Suriname toad, have evolved other ways of getting around this limitation). The first amniotes apparently arose in the middle Carboniferous from the ancestral reptiliomorphs. Within a few million years, two important amniote lineages became distinct: synapsids, from which mammals are descended, and sauropsids, from which lizards, snakes, turtles/tortoises, crocodilians, dinosaurs, and birds are descended. The earliest known fossils of synapsids and sauropsids (such as Archaeothyris and Hylonomus, respectively) date from about 320 to 315 million years ago. The times of origin are difficult to know, because vertebrate fossils from the late Carboniferous are very rare, and therefore the actual first occurrences of each of these types of animal might have been considerably earlier than the first fossil. Synapsids Synapsid skulls are identified by the distinctive pattern of the holes behind each eye, which served the following purposes: made the skull lighter without sacrificing strength. saved energy by using less bone. probably provided attachment points for jaw muscles. Having attachment points further away from the jaw made it possible for the muscles to be longer and therefore to exert a strong pull over a wide range of jaw movement without being stretched or contracted beyond their optimum range. A number of creatures often – and incorrectly – believed to be dinosaurs, hence part of the reptile lineage and sauropsids, were in fact synapsids. This includes the well-known Dimetrodon. Terms used for discussing non-mammalian synapsids When referring to the ancestors and close relatives of mammals, paleontologists also use the following terms of convenience: Pelycosaurs — all synapsids, and all of their descendants, except for therapsids – the eventual ancestor of mammals. The pelycosaurs included the largest land vertebrates of the Early Permian, such as the 6 metre (20 foot)-long Cotylorhynchus hancocki. Among the other large pelycosaurs were Dimetrodon grandis and Edaphosaurus cruciger. Stem mammals (sometimes called protomammals or paramammals, and previously called mammal-like reptiles) — all synapsids, and all of their descendants, except for mammals themselves. Stem mammals therefore include all pelycosaurs, and also all non-mammalian therapsids. Traditionally these were known as "mammal-like reptiles", but this is incorrect; terms such as "stem mammal" are preferred instead, because these synapsids were neither reptiles nor even part of reptile lineage. Therapsids Therapsids descended from sphenacodonts, a primitive synapsid, in the middle Permian, and took over from them as the dominant land vertebrates. They differ from earlier synapsids in several features of the skull and jaws, including larger temporal fenestrae and incisors that are equal in size. The therapsid lineage then went through several stages, leading to the evolution of cynodonts in the late Permian, some of which had begun to resemble early mammals: gradual development of a bony secondary palate. Most books and articles interpret this as a prerequisite for the evolution of mammals' high metabolic rate, because it enabled these animals to eat and breathe at the same time. But some scientists point out that some modern ectotherms use a fleshy secondary palate to separate the mouth from the airway, and that a bony palate provides a surface on which the tongue can manipulate food, facilitating chewing rather than breathing. The interpretation of the bony secondary palate as an aid to chewing also suggests the development of a faster metabolism, because chewing reduces the size of food particles delivered to the stomach and can therefore speed their digestion. In mammals, the palate is formed by two specific bones, but various Permian therapsids had other combinations of bones in the right places to function as a palate. the dentary gradually becomes the main bone of the lower jaw. Therapsid family tree A simplified phylogenetic tree showing only what is most relevant to the evolution of mammals is shown below: Only the dicynodonts, therocephalians, and cynodonts survived into the Triassic. Biarmosuchia The Biarmosuchia were the most primitive and pelycosaur-like of the therapsids. Dinocephalians Dinocephalians ("terrible heads") included both carnivores and herbivores. They were large; Anteosaurus was up to long. Some of the carnivores had semi-erect hindlimbs, but all dinocephalians had sprawling forelimbs. In many ways they were very primitive therapsids; for example, they had no secondary palate and their jaws were rather "reptilian". Anomodonts The anomodonts ("anomalous teeth") were among the most successful of the herbivorous therapsids — one sub-group, the dicynodonts, survived to the end of the Triassic. But dicynodonts were very different from modern herbivorous mammals, as their only teeth were a pair of fangs in the upper jaw (lost in some derived kannemeyeriiformes) and it is generally agreed that they had beaks like those of birds or ceratopsians. Theriodonts The theriodonts ("beast teeth") and their descendants had jaw joints in which the articular bone of the lower jaw tightly gripped the very small quadrate bone of the skull. This allowed a much wider gape and allowed one group, the carnivorous gorgonopsians ("gorgon faces"), to develop "sabre teeth". However, the jaw hinge of the theriodont had a longer term significance — the much reduced size of the quadrate bone was an important step in the development of the mammalian jaw joint and middle ear. The gorgonopsians still had some primitive features: no bony secondary palate (other bones in the right places perform the same functions); sprawling forelimbs; hindlimbs that could operate in both sprawling and erect postures. The therocephalians ("beast heads"), which appear to have arisen at about the same time as the gorgonopsians, had additional mammal-like features, e.g. their finger and toe bones had the same number of phalanges (segments) as in early mammals (and the same number that primates have, including humans). Numerous Changhsingian coprolites that possibly belong to therocephalians and indeterminate basal archosaurs (proterosuchids) contain elongated hollow structures that could be remains of hair. That means therapsids were covered in hair as early as 252 million years ago. Cynodonts The cynodonts, a theriodont group that also arose in the late Permian, include the ancestors of all mammals. Cynodonts' mammal-like features include further reduction in the number of bones in the lower jaw, a secondary bony palate, cheek teeth with a complex pattern in the crowns, and a brain which filled the endocranial cavity. Multi-chambered burrows have been found, containing as many as 20 skeletons of the Early Triassic cynodont Trirachodon; the animals are thought to have been drowned by a flash flood. The extensive shared burrows indicate that these animals were capable of complex social behaviors. Their primitive synapsid and therapsid ancestors were very large (between ) but cynodonts gradually decreased in size (to ) even before the Permian-Triassic extinction event, probably due to competition with other therapsids. After the extinction event, the probainognathian cynodont group rapidly decreased in size (to ) due to new competition with archosaurs and transitioned to nocturnality, evolving nocturnal features, pulmonary alveoli, bronchioles and a developed diaphragm for a larger surface area for breathing, enucleated erythrocytes, a large intestine which bears a true colon after the cecum, endothermy, a hairy, glandular and thermoregulatory skin (which releases sebum and sweat), and a 4-chambered heart to maintain their high metabolism, larger brains, and fully upright hindlimb (forelimbs remained semi sprawling, and became like that only later, in therians). Some skin glands may have evolved into mammary glands in females for fulfilling the metabolic demands of their offspring (which increased 10 times). Many skeletal changes occurred also: the dentary bone became stronger and held differentiated teeth, for example, and the pair of nasal openings in the skull became fused. These evolutionary changes led to the first mammals (size around ). They appear to have evolved rapid growth and short lifespan, a life history trait also found in numerous modern small-bodied mammals. They also adapted to a burrowing lifestyle, losing their large tail-based leg muscles which allowed dinosaurs to become bipedal, which may explain why bipedal mammals are so rare. Triassic takeover The catastrophic mass extinction at the end of the Permian, around 252 million years ago, killed off about 70% of terrestrial vertebrate species and the majority of land plants. As a result, ecosystems and food chains collapsed, and the establishment of new stable ecosystems took about 30 million years. With the disappearance of the gorgonopsians, which were dominant predators in the late Permian, the cynodonts' principal competitors for dominance of the carnivorous niches were a previously obscure sauropsid group, the archosaurs, which includes the ancestors of crocodilians and dinosaurs. The archosaurs quickly became the dominant carnivores, a development often called the "Triassic takeover". Their success may have been due to the fact that the early Triassic was predominantly arid and therefore archosaurs' superior water conservation gave them a decisive advantage. All known archosaurs have glandless skins and eliminate nitrogenous waste in a uric acid paste containing little water, while the cynodonts probably excreted most such waste in a solution of urea, as mammals do today; considerable water is required to keep urea dissolved. However, this theory has been questioned, since it implies synapsids were necessarily less advantaged in water retention, that synapsid decline coincides with climate changes or archosaur diversity (neither of which has been tested) and the fact that desert-dwelling mammals are as well adapted in this department as archosaurs, and some cynodonts like Trucidocynodon were large-sized predators. The Triassic takeover was probably a vital factor in the evolution of the mammals. Two groups stemming from the early cynodonts were successful in niches that had minimal competition from the archosaurs: the tritylodonts, which were herbivores, and the mammals, most of which were small nocturnal insectivores (although some, like Sinoconodon, were carnivores that fed on vertebrate prey, while still others were herbivores or omnivores). As a result: The therapsid trend towards differentiated teeth with precise occlusion accelerated, because of the need to hold captured arthropods and crush their exoskeletons. As the body length of the mammals' ancestors fell below , advances in thermal insulation and temperature regulation would have become necessary for nocturnal life. Acute senses of hearing and smell became vital. This accelerated the development of the mammalian middle ear (though the complete detachment of the middle ear bones from the jaw happened independently in monotremes). The increase in the size of the olfactory lobes of the brain increased brain weight as a percentage of total body weight. Brain tissue requires a disproportionate amount of energy. The need for more food to support the enlarged brains increased the pressures for improvements in insulation, temperature regulation and feeding. Probably as a side-effect of the nocturnal life, mammals lost two of the four cone opsins, photoreceptors in the retina, present in the eyes of the earliest amniotes. Paradoxically, this might have improved their ability to discriminate colors in dim light. This retreat to a nocturnal role is called a nocturnal bottleneck, and is thought to explain many of the features of mammals. From cynodonts to crown mammals Fossil record Mesozoic synapsids that had evolved to the point of having a jaw joint composed of the dentary and squamosal bones are preserved in few good fossils, mainly because they were mostly smaller than rats: They were largely restricted to environments that are less likely to provide good fossils. Floodplains as the best terrestrial environments for fossilization provide few mammal fossils, because they are dominated by medium to large animals, and the mammals could not compete with archosaurs in the medium to large size range. Their delicate bones were vulnerable to being destroyed before they could be fossilized — by scavengers (including fungi and bacteria) and by being trodden on. Small fossils are harder to spot and more vulnerable to being destroyed by weathering and other natural stresses before they are discovered. In the past years, however, the number of Mesozoic fossil mammals has increased decisively; only 116 genera were known in 1979, for example, but about 310 in 2007, with an increase in quality such that "at least 18 Mesozoic mammals are represented by nearly complete skeletons". Mammals or mammaliaforms Some writers restrict the term "mammal" to the crown group mammals, the group consisting of the most recent common ancestor of the monotremes, marsupials, and placentals, together with all the descendants of that ancestor. In an influential 1988 paper, Timothy Rowe advocated this restriction, arguing that "ancestry... provides the only means of properly defining taxa" and, in particular, that the divergence of the monotremes from the animals more closely related to marsupials and placentals "is of central interest to any study of Mammalia as a whole." To accommodate some related taxa falling outside the crown group, he defined the Mammaliaformes as comprising "the last common ancestor of Morganucodontidae and Mammalia [as he had defined the latter term] and all its descendants." Besides Morganucodontidae, the newly defined taxon includes Docodonta and Kuehneotheriidae. Though haramiyids have been referred to the mammals since the 1860s, Rowe excluded them from the Mammaliaformes as falling outside his definition, putting them in a larger clade, the Mammaliamorpha. Some writers have adopted this terminology noting, to avoid misunderstanding, that they have done so. Most paleontologists, however, still think that animals with the dentary-squamosal jaw joint and the sort of molars characteristic of modern mammals should formally be members of Mammalia. Where the ambiguity in the term "mammal" may be confusing, this article uses "mammaliaform" and "crown mammal". Family tree – cynodonts to crown group mammals (based on Cynodontia:Dendrogram – Palaeos) Morganucodontidae The Morganucodontidae first appeared in the late Triassic, about 205 million years ago. They are an excellent example of transitional fossils, since they have both the dentary-squamosal and articular-quadrate jaw joints. They were also one of the first discovered and most thoroughly studied of the mammaliaforms outside of the crown-group mammals, since an unusually large number of morganucodont fossils have been found. Docodonts Docodonts, among the most common Jurassic mammaliaforms, are noted for the sophistication of their molars. They are thought to have had general semi-aquatic tendencies, with the fish-eating Castorocauda ("beaver tail"), which lived in the mid-Jurassic about 164M years ago and was first discovered in 2004 and described in 2006, being the most well-understood example. Castorocauda was not a crown group mammal, but it is extremely important in the study of the evolution of mammals because the first find was an almost complete skeleton (a real luxury in paleontology) and it breaks the "small nocturnal insectivore" stereotype: It was noticeably larger than most Mesozoic mammaliaform fossils — about from its nose to the tip of its tail, and may have weighed . It provides the earliest absolutely certain evidence of hair and fur. Previously the earliest was Eomaia, a crown group mammal from about 125M years ago. It had aquatic adaptations including flattened tail bones and remnants of soft tissue between the toes of the back feet, suggesting that they were webbed. Previously the earliest known semi-aquatic mammaliaforms were from the Eocene, about 110M years later. Castorocauda'''s powerful forelimbs look adapted for digging. This feature and the spurs on its ankles make it resemble the platypus, which also swims and digs. Its teeth look adapted for eating fish: the first two molars had cusps in a straight row, which made them more suitable for gripping and slicing than for grinding; and these molars are curved backwards, to help in grasping slippery prey. Hadrocodium The family tree above shows Hadrocodium as a close relative of crown-group mammals. This mammaliaform, dated about 195 million years ago in the very early Jurassic, exhibits some important features: The jaw joint consists only of the squamosal and dentary bones, and the jaw contains no smaller bones to the rear of the dentary, unlike the therapsid design. In therapsids and early mammaliaforms, the eardrum may have stretched over a trough at the rear of the lower jaw. But Hadrocodium had no such trough, which suggests its ear was part of the cranium, as it is in crown-group mammals — and hence that the former articular and quadrate had migrated to the middle ear and become the malleus and incus. On the other hand, the dentary has a "bay" at the rear that mammals lack. This suggests that Hadrocodium's dentary bone retained the same shape that it would have had if the articular and quadrate had remained part of the jaw joint, and therefore that Hadrocodium or a very close ancestor may have been the first to have a fully mammalian middle ear. Therapsids and earlier mammaliaforms had their jaw joints very far back in the skull, partly because the ear was at the rear end of the jaw but also had to be close to the brain. This arrangement limited the size of the braincase, because it forced the jaw muscles to run round and over it. Hadrocodium's braincase and jaws were no longer bound to each other by the need to support the ear, and its jaw joint was further forward. In its descendants or those of animals with a similar arrangement, the brain case was free to expand without being constrained by the jaw and the jaw was free to change without being constrained by the need to keep the ear near the brain — in other words it now became possible for mammaliaforms both to develop large brains and to adapt their jaws and teeth in ways that were purely specialized for eating. Kuehneotheriidae The family Kuehneotheriidae, known from the Late Triassic and Early Jurassic, was originally classified as part of either 'Symmetrodonta' or 'Pantotheria' based on their tooth structure, with Kuehneotherium once being considered the oldest known representative of Theria. They have since been recovered as among the closest relatives of crown-group mammals. As only tooth fossils have been discovered, they nevertheless remain poorly known, and have rarely been included in phylogenetic studies. Earliest crown mammals The crown group mammals, sometimes called 'true mammals', are the extant mammals and their close relatives back to their last common ancestor. Since this group has living members, DNA analysis can be applied in an attempt to explain the evolution of features that do not appear in fossils. This endeavor often involves molecular phylogenetics, a technique that has become popular since the mid-1980s. Family tree of early crown mammals Cladogram after Z.-X Luo († marks extinct groups) and Hackländer. Color vision Early amniotes had four opsins in the cones of their retinas to use for distinguishing colours: one sensitive to red, one to green, and two corresponding to different shades of blue. The green opsin was not inherited by any crown mammals, but all normal individuals did inherit the red one. Early crown mammals thus had three cone opsins, the red one and both of the blues. All their extant descendants have lost one of the blue-sensitive opsins but not always the same one: monotremes retain one blue-sensitive opsin, while marsupials and placentals retain the other (except cetaceans, which later lost the other blue opsin as well). Some placentals and marsupials, including higher primates, subsequently evolved green-sensitive opsins; like early crown mammals, therefore, their vision is trichromatic. Australosphenida and Ausktribosphenidae Ausktribosphenidae is a group name that has been given to some rather puzzling finds that: appear to have tribosphenic molars, a type of tooth that is otherwise known only in placentals and marsupials. come from mid-Cretaceous deposits in Australia — but Australia was connected only to Antarctica, and placentals originated in the Northern Hemisphere and were confined to it until continental drift formed land connections from North America to South America, from Asia to Africa and from Asia to India. are represented only by teeth and jaw fragments, which is not very helpful. Australosphenida is a group that has been defined in order to include the Ausktribosphenidae and monotremes. Asfaltomylos (mid- to late Jurassic, from Patagonia) has been interpreted as a basal australosphenid (animal that has features shared with both Ausktribosphenidae and monotremes; lacks features that are peculiar to Ausktribosphenidae or monotremes; also lacks features that are absent in Ausktribosphenidae and monotremes) and as showing that australosphenids were widespread throughout Gondwanaland (the old Southern Hemisphere super-continent). Recent analysis of Teinolophos, which lived somewhere between 121 and 112.5 million years ago, suggests that it was a "crown group" (advanced and relatively specialised) monotreme. This was taken as evidence that the basal (most primitive) monotremes must have appeared considerably earlier, but this has been disputed (see the following section). The study also indicated that some alleged Australosphenids were also "crown group" monotremes (e.g. Steropodon) and that other alleged Australosphenids (e.g. Ausktribosphenos, Bishops, Ambondro, Asfaltomylos) are more closely related to and possibly members of the Therian mammals (group that includes marsupials and placentals, see below). MonotremesTeinolophos, from Australia, is the earliest known monotreme. A 2007 study (published 2008) suggests that it was not a basal (primitive, ancestral) monotreme but a full-fledged platypus, and therefore that the platypus and echidna lineages diverged considerably earlier. A more recent study (2009), however, has suggested that, while Teinolophos was a type of platypus, it was also a basal monotreme and predated the radiation of modern monotremes. The semi-aquatic lifestyle of platypuses prevented them from being outcompeted by the marsupials that migrated to Australia millions of years ago, since joeys need to remain attached to their mothers and would drown if their mothers ventured into water (though there are exceptions like the water opossum and the lutrine opossum; however, they both live in South America and thus do not come into contact with monotremes). Genetic evidence has determined that echidnas diverged from the platypus lineage as recently as 19-48M, when they made their transition from semi-aquatic to terrestrial lifestyle. Monotremes have some features that may be inherited from the cynodont ancestors: like lizards and birds, they use the same orifice to urinate, defecate and reproduce ("monotreme" means "one hole"). they lay eggs that are leathery and uncalcified, like those of lizards, turtles and crocodilians. Unlike other mammals, female monotremes do not have nipples and feed their young by "sweating" milk from patches on their bellies. These features are not visible in fossils, and the main characteristics from paleontologists' point of view are: a slender dentary bone in which the coronoid process is small or non-existent. the external opening of the ear lies at the posterior base of the jaw. the jugal bone is small or non-existent. a primitive pectoral girdle with strong ventral elements: coracoids, clavicles and interclavicle. Note: therian mammals have no interclavicle. sprawling or semi-sprawling forelimbs. Multituberculates Multituberculates (named for the multiple tubercles on their "molars") are often called the "rodents of the Mesozoic", but this is an example of convergent evolution rather than meaning that they are closely related to the Rodentia. They existed for approximately 120 million years—the longest fossil history of any mammal lineage—but were eventually outcompeted by rodents, becoming extinct during the early Oligocene. Some authors have challenged the phylogeny represented by the cladogram above. They exclude the multituberculates from the mammalian crown group, holding that multituberculates are more distantly related to extant mammals than even the Morganucodontidae. Multituberculates are like undisputed crown mammals in that their jaw joints consist of only the dentary and squamosal bones-whereas the quadrate and articular bones are part of the middle ear; their teeth are differentiated, occlude, and have mammal-like cusps; they have a zygomatic arch; and the structure of the pelvis suggests that they gave birth to tiny helpless young, like modern marsupials. On the other hand, they differ from modern mammals: Their "molars" have two parallel rows of tubercles, unlike the tribosphenic (three-peaked) molars of uncontested early crown mammals. The chewing action differs in that undisputed crown mammals chew with a side-to-side grinding action, which means that the molars usually occlude on only one side at a time, while multituberculates' jaws were incapable of side-to-side movement—they chewed, rather, by dragging the lower teeth backwards against the upper ones as the jaw closed. The anterior (forward) part of the zygomatic arch mostly consists of the maxilla (upper jawbone) rather than the jugal, a small bone in a little slot in the maxillary process (extension). The squamosal does not form part of the braincase. The rostrum (snout) is unlike that of undisputed crown mammals; in fact it looks more like that of a pelycosaur, such as Dimetrodon. The multituberculate rostrum is box-like, with the large flat maxillae forming the sides, the nasal the top, and the tall premaxilla at the front. Theria Theria ("beasts") is the clade originating with the last common ancestor of the Eutheria (including placentals) and Metatheria (including marsupials). Common features include: no interclavicle. coracoid bones non-existent or fused with the shoulder blades to form coracoid processes. a type of crurotarsal ankle joint in which: the main joint is between the tibia and astragalus; the calcaneum has no contact with the tibia but forms a heel to which muscles can attach. (The other well-known type of crurotarsal ankle is seen in crocodilians and works differently — most of the bending at the ankle is between the calcaneum and astragalus). tribosphenic molars. Metatheria The living Metatheria are all marsupials (animals with pouches). A few fossil genera, such as the Mongolian late Cretaceous Asiatherium, may be marsupials or members of some other metatherian group(s). The oldest known metatherian is Sinodelphys, found in 125M-year-old early Cretaceous shale in China's northeastern Liaoning Province. The fossil is nearly complete and includes tufts of fur and imprints of soft tissues. Didelphimorphia (common opossums of the Western Hemisphere) first appeared in the late Cretaceous and still have living representatives, probably because they are mostly semi-arboreal unspecialized omnivores. Tracks from the Early Cretaceous of Angola show the existence of raccoon-size mammals 118 million years ago. The best-known feature of marsupials is their method of reproduction: The mother develops a kind of yolk sack in her womb that delivers nutrients to the embryo. Embryos of bandicoots, koalas and wombats additionally form placenta-like organs that connect them to the uterine wall, although the placenta-like organs are smaller than in placental mammals and it is not certain that they transfer nutrients from the mother to the embryo. Pregnancy is very short, typically four to five weeks. The embryo is born at a very early stage of development, and is usually less than long at birth. It has been suggested that the short pregnancy is necessary to reduce the risk that the mother's immune system will attack the embryo. The newborn marsupial uses its forelimbs (with relatively strong hands) to climb to a nipple, which is usually in a pouch on the mother's belly. The mother feeds the baby by contracting muscles over her mammary glands, as the baby is too weak to suck. The newborn marsupial's need to use its forelimbs in climbing to the nipple was historically thought to have restricted metatherian evolution, as it was assumed that the forelimb could not become specialised intro structures like wings, hooves or flippers. However, several bandicoots, most notably the pig-footed bandicoot, have true hooves similar to those of placental ungulates, and several marsupial gliders have evolved. Although some marsupials look very like some placentals (the thylacine, "marsupial tiger" or "marsupial wolf" is a good example), marsupial skeletons have some features that distinguish them from placentals: Some, including the thylacine, have four molars; whereas no known placental has more than three. All have a pair of palatal fenestrae, window-like openings on the bottom of the skull (in addition to the smaller nostril openings). Marsupials also have a pair of marsupial bones (sometimes called "epipubic bones"), which support the pouch in females. But these are not unique to marsupials, since they have been found in fossils of multituberculates, monotremes, and even eutherians — so they are probably a common ancestral feature that disappeared at some point after the ancestry of living placental mammals diverged from that of marsupials. Some researchers think the epipubic bones' original function was to assist locomotion by supporting some of the muscles that pull the thigh forwards. Eutheria The time of appearance of the earliest eutherians has been a matter of controversy. On one hand, recently discovered fossils of Juramaia have been dated to 160 million years ago and classified as eutherian. Fossils of Eomaia from 125 million years ago in the Early Cretaceous have also been classified as eutherian. A recent analysis of phenomic characters, however, classified Eomaia as pre-eutherian and reported that the earliest clearly eutherian specimens came from Maelestes, dated to 91 million years ago. That study also reported that eutherians did not significantly diversify until after the catastrophic extinction at the Cretaceous–Paleogene boundary, about 66 million years ago.Eomaia was found to have some features that are more like those of marsupials and earlier metatherians: Epipubic bones extending forwards from the pelvis, which are not found in any modern placental, but are found in all other mammals — early mammaliaforms, non-placental eutherians, marsupials, and monotremes — as well as in the cynodont therapsids that are closest to mammals. Their function is to stiffen the body during locomotion. This stiffening would be harmful in pregnant placentals, whose abdomens need to expand. A narrow pelvic outlet, which indicates that the young were very small at birth and therefore pregnancy was short, as in modern marsupials. This suggests that the placenta was a later development. Five incisors in each side of the upper jaw. This number is typical of metatherians, and the maximum number in modern placentals is three, except for homodonts, such as the armadillo. But Eomaia's molar to premolar ratio (it has more pre-molars than molars) is typical of eutherians, including placentals, and not normal in marsupials.Eomaia also has a Meckelian groove, a primitive feature of the lower jaw that is not found in modern placental mammals. These intermediate features are consistent with molecular phylogenetics estimates that the placentals diversified about 110M years ago, 15M years after the date of the Eomaia fossil.Eomaia also has many features that strongly suggest it was a climber, including several features of the feet and toes; well-developed attachment points for muscles that are used a lot in climbing; and a tail that is twice as long as the rest of the spine. Placentals' best-known feature is their method of reproduction: The embryo attaches itself to the uterus via a large placenta via which the mother supplies food and oxygen and removes waste products. Pregnancy is relatively long and the young are fairly well developed at birth. In some species (especially herbivores living on plains) the young can walk and even run within an hour of birth. It has been suggested that the evolution of placental reproduction was made possible by retroviruses that: make the interface between the placenta and uterus into a syncytium, i.e. a thin layer of cells with a shared external membrane. This allows the passage of oxygen, nutrients and waste products, but prevents the passage of blood and other cells that would cause the mother's immune system to attack the fetus. reduce the aggressiveness of the mother's immune system, which is good for the foetus but makes the mother more vulnerable to infections. From a paleontologist's point of view, eutherians are mainly distinguished by various features of their teeth, ankles and feet. Expansion of ecological niches in the Mesozoic thumb\|Skull cast of Late Cretaceous Didelphodon, showing its robust teeth adapted to a durophagous diet. Generally speaking, most species of mammaliaforms did occupy the niche of small, nocturnal insectivores, but recent finds, mainly in China, show that some species and especially crown group mammals were larger and that there was a larger variety of lifestyles than previously thought. For example: The therian Patagomaia, found in the Late Cretaceous Chorrillo Formation (Argentina) is the largest known Mesozoic mammal, weighing an estimated . Adalatherium hui is a large sized, erect limbed herbivore from the Cretaceous of Madagascar.Castorocauda, a member of Docodonta which lived in the middle Jurassic about 164 million years, was about long, weighed , had a beaver-like tail that was adapted for swimming, limbs adapted for swimming and digging, and teeth adapted for eating fish. Another docodont, Haldanodon, also had semi-aquatic habits, and indeed aquatic tendencies were probably common among docodonts based on their prevalence in wetland environments. The eutriconodonts Liaoconodon and Yanoconodon have more recently also have been suggested to be freshwater swimmers, lacking Castorocaudas powerful tail but possessing paddle-like limbs; the eutriconodont Astroconodon has similarly been suggested as being semi-aquatic in the past, albeit to less convincing evidence. Multituberculates are allotherians that survived for over 125 million years (from mid-Jurassic, about 160M years ago, to late Eocene, about 35M years ago) are often called the "rodents of the Mesozoic". As noted above, they may have given birth to tiny live neonates rather than laying eggs. Fruitafossor, from the late Jurassic period about 150 million years ago, was about the size of a chipmunk and its teeth, forelimbs and back suggest that it broke open the nest of social insects to prey on them (probably termites, as ants had not yet appeared). Similarly, the gobiconodontid Spinolestes possessed adaptations for fossoriality and convergent traits with placental xenarthrans like scutes and xenarthrous vertebrae, so it too might have had anteater like habits. It is also notable for the presence of quills akin to those of modern spiny mice. Volaticotherium, from the boundary the early Cretaceous about 125M years ago, is the earliest-known gliding mammal and had a gliding membrane that stretched out between its limbs, rather like that of a modern flying squirrel. This also suggests it was active mainly during the day. The closely related Argentoconodon also shows similar adaptations that may also suggest aerial locomotion. Repenomamus, a eutriconodont from the early Cretaceous 130 million years ago, was a stocky, badger-like predator that sometimes preyed on young dinosaurs. Two species have been recognized, one more than long and weighing about , the other less than long and weighing . Schowalteria is a Late Cretaceous species almost as large if not larger than R. giganticus that shows speciations towards herbivory, comparable to those of modern ungulates. Zhelestidae is a lineage of Late Cretaceous herbivorous eutherians, to the point of being mistaken for stem-ungulates. Similarly, mesungulatids are also fairly large sized herbivorous mammals from the Late Cretaceous Deltatheroidans were metatherians that were specialised towards carnivorous habits, and possible forms like Oxlestes and Khudulestes might have been among the largest Mesozoic mammals, though their status as deltatheroidans is questionable. Ichthyoconodon, a eutriconodont from the Berriasian of Morocco, is currently known from molariforms found in marine deposits. These teeth are sharp-cusped and similar in shape to those of piscivorous mammals, and unlike the teeth of contemporary mammals they do not show degradation, so rather than being carried down by river deposits the animal died in situ or close. This has been taken to mean that it was a marine mammal, likely one of the few examples known from the Mesozoic. Didelphodon is a Late Cretaceous riverine species of stagodontid marsupialiform with a durophagous dentition, robust jaws similar to a modern Tasmanian devil, and a postcranial skeleton very similar in size and shape to an otter. This animal has been lauded as the strongest bite of all Mesozoic mammals. It possibly specialized on eating freshwater crabs and molluscs. Tracks of a raccoon-sized mammaliaform representing the morphofamily Ameghinichnidae are described from the Early Cretaceous (late Aptian) Calonda Formation (Angola) by Mateus et al. (2017), who name a new ichnotaxon Catocapes angolanus. A gobiconodontid was preserved attacking a substantially larger dinosaur. A study on Mesozoic mammaliaforms suggests that they were a primary factor in constraining mammalian body size, rather than solely competition from dinosaurs. In general, it appears mammal faunas on southern continents had attained larger body sizes than those of northern continents. Evolution of major groups of living mammals There are currently vigorous debates between traditional paleontologists and molecular phylogeneticists about how and when the modern groups of mammals diversified, especially the placentals. Generally, the traditional paleontologists date the appearance of a particular group by the earliest known fossil whose features make it likely to be a member of that group, while the molecular phylogeneticists suggest that each lineage diverged earlier (usually in the Cretaceous) and that the earliest members of each group were anatomically very similar to early members of other groups and differed only in their genetics. These debates extend to the definition of and relationships between the major groups of placentals. Molecular phylogenetics-based family tree of placental mammals Molecular phylogenetics uses features of organisms' genes to work out family trees in much the same way as paleontologists do with features of fossils — if two organisms' genes are more similar to each other than to those of a third organism, the two organisms are more closely related to each other than to the third. Molecular phylogeneticists have proposed a family tree that is both broadly similar to but has notable differences from that of the paleontologists. Like paleontologists, molecular phylogeneticists have differing ideas about various details, but here is a typical family tree according to molecular phylogenetics: Note that the diagram shown here omits extinct groups, as one cannot extract DNA from most fossils. Some finer-level subdivisions are glossed-over. Here are the most significant of the differences between this family tree and the one familiar to paleontologists: The top-level division is between Atlantogenata and Boreoeutheria, instead of between Xenarthra and the rest. However, analysis of transposable element insertions supports a three-way top-level split between Xenarthra, Afrotheria and Boreoeutheria and the Atlantogenata clade does not receive significant support in recent distance-based molecular phylogenetics. Afrotheria contains several groups that are only distantly related according to the paleontologists' version: Afroinsectiphilia ("African insectivores"), Tubulidentata (aardvarks, which paleontologists regard as much closer to odd-toed ungulates than to other members of Afrotheria), Macroscelidea (elephant shrews, usually regarded as close to rabbits and rodents). The only members of Afrotheria that paleontologists would regard as closely related are Hyracoidea (hyraxes), Proboscidea (elephants) and Sirenia (manatees, dugongs). Members of the defunct order of Insectivores are divided among three clades: One clade is part of Afrotheria, and the other two clades are distinct sub-groups within Boreoeutheria. Bats are closer to Carnivora and odd-toed ungulates than to Primates and Dermoptera (colugos). Perissodactyla (odd-toed ungulates) are closer to Carnivora and bats than to Artiodactyla (even-toed ungulates). The grouping together of the Afrotheria has some geological justification: All surviving members of the Afrotheria originate from South American or (mainly) African lineages — even the Indian elephant, which diverged from an African lineage about . As Pangaea broke up, Africa and South America separated from the other continents less than 150M years ago, and from each other between 100M and 80M years ago. So it would not be surprising if the earliest eutherian immigrants into Africa and South America were isolated there and radiated into all the available ecological niches. Nevertheless, these proposals have been controversial. Paleontologists naturally insist that fossil evidence must take priority over deductions from samples of the DNA of modern animals. More surprisingly, these new family trees have been criticised by other molecular phylogeneticists, sometimes quite harshly: Mitochondrial DNA's mutation rate in mammals varies from region to region — some parts hardly ever change and some change extremely quickly and even show large variations between individuals within the same species. Mammalian mitochondrial DNA mutates so fast that it causes a problem called "saturation", where random noise drowns out any information that may be present. If a particular piece of mitochondrial DNA mutates randomly every few million years, it will have changed several times in the 60 to 75M years since the major groups of placental mammals diverged. Timing of placental evolution Recent molecular phylogenetic studies suggest that most placental orders diverged late in the Cretaceous period, about 100 to 85 million years ago, but that modern families first appeared later, in the late Eocene and early Miocene epochs of the Cenozoic period. Fossil-based analyses, on the contrary, limit the placentals to the Cenozoic. Many Cretaceous fossil sites contain well-preserved lizards, salamanders, birds, and mammals, but not the modern forms of mammals. It is possible that they simply did not exist, and that the molecular clock runs fast during major evolutionary radiations. On the other hand, there is fossil evidence from of hoofed mammals that may be ancestors of modern ungulates. Fossils of the earliest members of most modern groups date from the Paleocene, a few date from later and very few from the Cretaceous, before the extinction of the dinosaurs. But some paleontologists, influenced by molecular phylogenetic studies, have used statistical methods to extrapolate backwards from fossils of members of modern groups and concluded that primates arose in the late Cretaceous. However, statistical studies of the fossil record confirm that mammals were restricted in size and diversity right to the end of the Cretaceous, and rapidly grew in size and diversity during the Early Paleocene. Evolution of mammalian features Jaws and middle ears Hadrocodium, whose fossils date from the early Jurassic, provides the first clear evidence of fully mammalian jaw joints and middle ears, in which the jaw joint is formed by the dentary and squamosal bones while the articular and quadrate move to the middle ear, where they are known as the incus and malleus. One analysis of the monotreme Teinolophos suggested that this animal had a pre-mammalian jaw joint formed by the angular and quadrate bones and that the definitive mammalian middle ear evolved twice independently, in monotremes and in therian mammals, but this idea has been disputed. In fact, two of the suggestion's authors co-authored a later paper that reinterpreted the same features as evidence that Teinolophos was a full-fledged platypus, which means it would have had a mammalian jaw joint and middle ear. Lactation It has been suggested that lactation's original function was to keep eggs moist. Much of the argument is based on monotremes (egg-laying mammals): While the amniote egg is usually described as able to evolve away from water, most reptile eggs actually need moisture if they are not to dry out. Monotremes do not have nipples, but secrete milk from a hairy patch on their bellies. During incubation, monotreme eggs are covered in a sticky substance whose origin is not known. Before the eggs are laid, their shells have only three layers. Afterwards, a fourth layer appears with a composition different from that of the original three. The sticky substance and the fourth layer may be produced by the mammary glands. If so, that may explain why the patches from which monotremes secrete milk are hairy. It is easier to spread moisture and other substances over the egg from a broad, hairy area than from a small, bare nipple. Later research demonstrated that caseins already appeared in the common mammalian ancestor approximately 200–310 million years ago. The question of whether secretions of a substance to keep eggs moist translated into actual lactation in therapsids is open. A small mammaliomorph called Sinocodon, generally assumed to be the sister group of all later mammals, had front teeth in even the smallest individuals. Combined with a poorly ossified jaw, they very probably did not suckle. Thus suckling may have evolved right at the pre-mammal/mammal transition. However, tritylodontids, generally assumed to be more basal, show evidence of suckling. Morganucodontans, also assumed to be basal Mammaliaformes, also show evidence of lactation. Digestive system The evolution of the digestive system has formed a significant influence in mammal evolution. With the emergence of mammals, the digestive system was modified in a variety of ways depending on the animal's diet. For example, cats and most carnivores have simple large intestines, while the horse as a herbivore has a voluminous large intestine. An ancestral feature of ruminants is their multi-chambered (usually four-chambered) stomach, which evolved about 50 million years ago. Along with morphology of the gut, gastric acidity has been proposed as a key factor shaping the diversity and composition of microbial communities found in the vertebrate gut. Comparisons of stomach acidity across trophic groups in mammal and bird taxa show that scavengers and carnivores have significantly higher stomach acidities compared to herbivores or carnivores feeding on phylogenetically distant prey such as insects or fish. Despite the lack of fossilization of the gut, microbial evolution of the gut can be inferred from the interrelationships of existing animals, microbes and probable foodstuffs. Mammals are metagenomic, in that they are composed of not only their own genes, but also those of all of their associated microbes. Gut microbiota has co-diversified as mammalian species have evolved. Recent studies indicate that adaptive divergence between mammalian species is shaped in part by changes in the gut microbiota. The house mouse may have evolved not only with, but also in response to, the unique bacteria inhabiting its gut. Hair and fur The first clear evidence of hair or fur is in fossils of Castorocauda and Megaconus, from 164M years ago in the mid-Jurassic. As both mammals Megaconus and Castorocauda have a double coat of hair, with both guard hairs and an undercoat, it may be assumed that their last common ancestor did as well. More recently, the discovery of hair remnants in Permian coprolites pushes back the origin of mammalian hair much further back in the synapsid line to Paleozoic therapsids. In the mid-1950s, some scientists interpreted the foramina (passages) in the maxillae (upper jaws) and premaxillae (small bones in front of the maxillae) of cynodonts as channels that supplied blood vessels and nerves to vibrissae (whiskers) and suggested that this was evidence of hair or fur. It was soon pointed out, however, that foramina do not necessarily show that an animal had vibrissae; the modern lizard Tupinambis has foramina that are almost identical to those found in the non-mammalian cynodont Thrinaxodon. Popular sources, nevertheless, continue to attribute whiskers to Thrinaxodon. A trace fossil from the Lower Triassic had been erroneously regarded as a cynodont footprint showing hair, but this interpretation has been refuted. A study of cranial openings for facial nerves connected whiskers in extant mammals indicate the Prozostrodontia, small immediate ancestors of mammals, presented whiskers similar to mammals, but that less advanced therapsids would either have immobile whiskers or no whisker at all. Fur may have evolved from whiskers. Whiskers themselves may have evolved as a response to nocturnal and/or burrowing lifestyle. Ruben & Jones (2000) note that the Harderian glands, which secrete lipids for coating the fur, were present in the earliest mammals like Morganucodon, but were absent in near-mammalian therapsids like Thrinaxodon. The Msx2 gene associated with hair follicle maintenance is also linked to the closure of the parietal eye in mammals, indicating that fur and lack of pineal eye is linked. The pineal eye is present in Thrinaxodon, but absent in more advanced cynognaths (the Probainognathia). Insulation is the "cheapest" way to maintain a fairly constant body temperature, without consuming energy to produce more body heat. Therefore, the possession of hair or fur would be good evidence of homeothermy, but would not be such strong evidence of a high metabolic rate. Erect limbs Understanding of the evolution of erect limbs in mammals is incomplete — living and fossil monotremes have sprawling limbs. Some scientists think that the parasagittal (non-sprawling) limb posture is limited to the Boreosphenida, a group that contains the therians but not, for example, the multituberculates. In particular, they attribute a parasagittal stance to the therians Sinodelphys and Eomaia, which means that the stance had arisen by 125 million years ago, in the Early Cretaceous. However, they also discuss that earlier mammals had more erect forelimbs as opposed to the more sprawling hindlimbs, a trend still continued to some extent in modern placentals and marsupials. Warm-bloodedness "Warm-bloodedness" is a complex and rather ambiguous term, because it includes some or all of the following:Endothermy, the ability to generate heat internally rather than via behaviors such as basking or muscular activity.Homeothermy, maintaining a fairly constant body temperature. Most enzymes have an optimum operating temperature; efficiency drops rapidly outside the preferred range. A homeothermic organism needs only to possess enzymes that function well in a small range of temperatures.Tachymetabolism', maintaining a high metabolic rate, particularly when at rest. This requires a fairly high and stable body temperature because of the Q10 effect: biochemical processes run about half as fast if an animal's temperature drops by 10 °C. Since scientists cannot know much about the internal mechanisms of extinct creatures, most discussion focuses on homeothermy and tachymetabolism. However, it is generally agreed that endothermy first evolved in non-mammalian synapsids such as dicynodonts, which possess body proportions associated with heat retention, high vascularised bones with Haversian canals, and possibly hair. More recently, it has been suggested that endothermy evolved as far back as Ophiacodon. Modern monotremes have a low body temperature compared to marsupials and placental mammals, around . Phylogenetic bracketing suggests that the body temperatures of early crown-group mammals were not less than that of extant monotremes. There is cytological evidence that the low metabolism of monotremes is a secondarily evolved trait. Respiratory turbinates Modern mammals have respiratory turbinates, convoluted structures of thin bone in the nasal cavity. These are lined with mucous membranes that warm and moisten inhaled air and extract heat and moisture from exhaled air. An animal with respiratory turbinates can maintain a high rate of breathing without the danger of drying its lungs out, and therefore may have a fast metabolism. Unfortunately these bones are very delicate and therefore have not yet been found in fossils. But rudimentary ridges like those that support respiratory turbinates have been found in advanced Triassic cynodonts, such as Thrinaxodon and Diademodon, which suggests that they may have had fairly high metabolic rates. Bony secondary palate Mammals have a secondary bony palate, which separates the respiratory passage from the mouth, allowing them to eat and breathe at the same time. Secondary bony palates have been found in the more advanced cynodonts and have been used as evidence of high metabolic rates. But some cold-blooded vertebrates have secondary bony palates (crocodilians and some lizards), while birds, which are warm-blooded, do not. Diaphragm A muscular diaphragm helps mammals to breathe, especially during strenuous activity. For a diaphragm to work, the ribs must not restrict the abdomen, so that expansion of the chest can be compensated for by reduction in the volume of the abdomen and vice versa. Diaphragms are known in caseid pelycosaurs, indicating an early origin within synapsids, though they were still fairly inefficient and likely required support from other muscle groups and limb motion. The advanced cynodonts have very mammal-like rib cages, with greatly reduced lumbar ribs. This suggests that these animals had more developed diaphragms, were capable of strenuous activity for fairly long periods and therefore had high metabolic rates. On the other hand, these mammal-like rib cages may have evolved to increase agility. However, the movement of even advanced therapsids was "like a wheelbarrow", with the hindlimbs providing all the thrust while the forelimbs only steered the animal, in other words advanced therapsids were not as agile as either modern mammals or the early dinosaurs. So the idea that the main function of these mammal-like rib cages was to increase agility is doubtful. Limb posture The therapsids had sprawling forelimbs and semi-erect hindlimbs. This suggests that Carrier's constraint would have made it rather difficult for them to move and breathe at the same time, but not as difficult as it is for animals such as lizards, which have completely sprawling limbs. Advanced therapsids may therefore have been significantly less active than modern mammals of similar size and so may have had slower metabolisms overall or else been bradymetabolic (lower metabolism when at rest). Brain Mammals are noted for their large brain size relative to body size, compared to other animal groups. Recent findings suggest that the first brain area to expand was that involved in smell. Scientists scanned the skulls of early mammal species dating back to 190–200 million years ago and compared the brain case shapes to earlier pre-mammal species; they found that the brain area involved in the sense of smell was the first to enlarge. This change may have allowed these early mammals to hunt insects at night when dinosaurs were not active. After the extinction of the dinosaurs 66 million years ago, mammals began to increase in body size as new niches became available, but their brain lagged behind their bodies for the first ten million years. Relative to body size the brain of Paleocene mammal was relatively smaller than that of Mesozoic mammals. It was not until the Eocene that the mammalian brains began to catch up with their bodies, particularly in certain areas associated with their senses. Testicular descent Sexual selection	Biology and health sciences	Basics_4	Biology
7096085	https://en.wikipedia.org/wiki/Construction%20aggregate	Construction aggregate	Construction aggregate, or simply aggregate, is a broad category of coarse- to medium-grained particulate material used in construction. Traditionally, it includes natural materials such as sand, gravel, crushed stone. As with other types of aggregates, it is a component of composite materials, particularly concrete and asphalt. Aggregates are the most mined materials in the world, being a significant part of 6 billion tons of concrete produced per year. Aggregate serves as reinforcement to add strength to the resulting material. Due to the relatively high hydraulic conductivity as compared to most soil types, aggregates are widely used in drainage applications such as foundation and French drains, septic drain fields, retaining wall drains, and roadside edge drains. Aggregates are also used as base material under building foundations, roads, and railroads (aggregate base). It has predictable, uniform properties, preventing differential settling under the road or building. Aggregates are also used as a low-cost extender that binds with more expensive cement or asphalt to form concrete. Although most kinds of aggregate require a form of binding agent, there are types of self-binding aggregate which require no form of binding agent. More recently, recycled concrete and geosynthetic materials have also been used as aggregates. Sources Sources for these basic materials can be grouped into three main areas: mining of mineral aggregate deposits, including sand, gravel, and stone; use of waste slag from the manufacture of iron and steel; and recycling of concrete, which is itself chiefly manufactured from mineral aggregates. In addition, there are some (minor) materials that are used as specialty lightweight aggregates: clay, pumice, perlite, and vermiculite. Other minerals include: basalt dolomite granite gravel limestone sand sandstone Specifications In Europe, sizing ranges are specified as d/D, where the d shows the smallest and D shows the largest square mesh grating that the particles can pass. Application-specific preferred sizings are covered in European Standard EN 13043 for road construction, EN 13383 for larger armour stone, EN 12620 for concrete aggregate, EN 13242 for base layers of road construction, and EN 13450 for railway ballast. The American Society for Testing and Materials publishes an exhaustive listing of specifications including ASTM D 692 and ASTM D 1073 for various construction aggregate products, which, by their individual design, are suitable for specific construction purposes. These products include specific types of coarse and fine aggregate designed for such uses as additives to asphalt and concrete mixes, as well as other construction uses. State transportation departments further refine aggregate material specifications in order to tailor aggregate use to the needs and available supply in their particular locations. History People have used sand and stone for foundations for thousands of years. Significant refinement of the production and use of aggregate occurred during the Roman Empire, which used aggregate to build its vast network of roads and aqueducts. The invention of concrete, which was essential to architecture utilizing arches, created an immediate, permanent demand for construction aggregates. Vitruvius writes in De architectura: Economy denotes the proper management of materials and of site, as well as a thrifty balancing of cost and common sense in the construction of works. This will be observed if, in the first place, the architect does not demand things which cannot be found or made ready without great expense. For example: it is not everywhere that there is plenty of pit-sand, rubble, fir, clear fir, and marble... Where there is no pit sand, we must use the kinds washed up by rivers or by the sea... and other problems we must solve in similar ways. Modern production The advent of modern blasting methods enabled the development of quarries, which are now used throughout the world, wherever competent bedrock deposits of aggregate quality exist. In many places, good limestone, granite, marble or other quality stone bedrock deposits do not exist. In these areas, natural sand and gravel are mined for use as aggregate. Where neither stone, nor sand and gravel, are available, construction demand is usually satisfied by shipping in aggregate by rail, barge or truck. Additionally, demand for aggregates can be partially satisfied through the use of slag and recycled concrete. However, the available tonnages and lesser quality of these materials prevent them from being a viable replacement for mined aggregates on a large scale. Large stone quarry and sand and gravel operations exist near virtually all population centers due to the high cost of transportation relative to the low value of the product. Trucking aggregate more than 40 kilometers is typically uneconomical. These are capital-intensive operations, utilizing large earth-moving equipment, belt conveyors, and machines specifically designed for crushing and separating various sizes of aggregate, to create distinct product stockpiles. According to the USGS, 2006 U.S. crushed stone production was 1.72 billion tonnes valued at $13.8 billion (compared to 1.69 billion tonnes valued at $12.1 billion in 2005), of which limestone was 1,080 million tonnes valued at $8.19 billion from 1,896 quarries, granite was 268 million tonnes valued at $2.59 billion from 378 quarries, trap rock was 148 million tonnes valued at $1.04 billion from 355 quarries, and the balance other kinds of stone from 729 quarries. Limestone and granite are also produced in large amounts as dimension stone. The great majority of crushed stone is moved by heavy truck from the quarry/plant to the first point of sale or use. According to the USGS, 2006 U.S. sand and gravel production was 1.32 billion tonnes valued at $8.54 billion (compared to 1.27 billion tonnes valued at $7.46 billion in 2005), of which 264 million tonnes valued at $1.92 billion was used as concrete aggregates. The great majority of this was again moved by truck, instead of by electric train. Currently, total U.S. aggregate demand by final market sector was 30%–35% for non-residential building (offices, hotels, stores, manufacturing plants, government and institutional buildings, and others), 25% for highways, and 25% for housing. Recycled materials Recycled material such as blast furnace and steel furnace slag can be used as aggregate or partly substitute for portland cement. Blast furnace and steel slag is either air-cooled or water-cooled. Air-cooled slag can be used as aggregate. Water-cooled slag produces sand-sized glass-like particles (granulated). Adding free lime to the water during cooling gives granulated slag hydraulic cementitious properties. In 2006, according to the USGS, air-cooled blast furnace slag sold or used in the U.S. was 7.3 million tonnes valued at $49 million, granulated blast furnace slag sold or used in the U.S. was 4.2 million tonnes valued at $318 million, and steel furnace slag sold or used in the U.S. was 8.7 million tonnes valued at $40 million. Air-cooled blast furnace slag sales in 2006 were for use in road bases and surfaces (41%), asphaltic concrete (13%), ready-mixed concrete (16%), and the balance for other uses. Granulated blast furnace slag sales in 2006 were for use in cementitious materials (94%), and the balance for other uses. Steel furnace slag sales in 2006 were for use in road bases and surfaces (51%), asphaltic concrete (12%), for fill (18%), and the balance for other uses. Recycled glass aggregate crushed to a small size is substituted for many construction and utility projects in place of pea gravel or crushed rock. Glass aggregate is not dangerous to handle. It can be used as pipe bedding—placed around sewer, storm water or drinking water pipes to transfer weight from the surface and protect the pipe. Another common use is as fill to bring the level of a concrete floor even with a foundation. Use of glass aggregate helps close the loop in glass recycling in many places where glass cannot be smelted into new glass. Aggregates themselves can be recycled as aggregates. Recyclable aggregate tends to be concentrated in urban areas. The supply of recycled aggregate depends on physical decay and demolition of structures. Mobile recycling plants eliminate the cost of transporting the material to a central site. The recycled material is typically of variable quality. Many aggregate products are recycled for other industrial purposes. Contractors save on disposal costs and less aggregate is buried or piled and abandoned. In Bay City, Michigan, for example, a recycle program exists for unused products such as mixed concrete, block, brick, gravel, pea stone, and other used materials. The material is crushed to provide subbase for roads and driveways, among other purposes. According to the USGS in 2006, 2.9 million tonnes of Portland cement concrete (including aggregate) worth $21.9 million was recycled, and 1.6 million tonnes of asphalt concrete (including aggregate) worth $11.8 million was recycled, both by crushed stone operations. Much more of both materials are recycled by construction and demolition firms not included in the USGS survey. For sand and gravel, the survey showed that 4.7 million tonnes of cement concrete valued at $32.0 million was recycled, and 6.17 million tonnes of asphalt concrete valued at $45.1 million was recycled. Again, more of both materials are recycled by construction and demolition firms not in this USGS survey. The Construction Materials Recycling Association indicates that there are 325 million tonnes of recoverable construction and demolition materials produced annually. Organic materials Many geosynthetic aggregates are made from recycled materials. Recyclable plastics can be reused in aggregates. For example, Ring Industrial Group's EZflow product lines are produced with geosynthetic aggregate pieces that are more than 99.9% recycled polystyrene. This polystyrene, otherwise destined for a landfill, is gathered, melted, mixed, reformulated and expanded to create low density aggregates that maintain high strength properties under compressive loads. Such geosynthetic aggregates replace conventional gravel while simultaneously increasing porosity, increasing hydraulic conductivity and eliminating the fine dust "fines" inherent to gravel aggregates which otherwise serve to clog and disrupt the operation of many drainage applications. Several groups have attempted to use minced tires as part of concrete aggregate. The result is tougher than regular concrete, because it can bend instead of breaking under pressure. However, tires reduce compressive strength partially because the cement bonds poorly with the rubber. Pores in the rubber fill with water when the concrete is mixed, but become voids as the concrete sets. One group put the concrete under pressure as it sets, reducing pore volumes. Recycled aggregates in the UK Recycled aggregate in the UK results from the processing of construction material. To ensure the aggregate is inert, it is manufactured from material tested and characterised under European Waste Codes. In 2008, 210 million tonnes of aggregate were produced including 67 million tonnes of recycled product, according to the Quarry Products Association. The Waste and Resource Action Programme has produced a Quality Protocol for the regulated production of recycled aggregates.	Technology	Building materials	null
1495085	https://en.wikipedia.org/wiki/Archaeopteris	Archaeopteris	Archaeopteris is an extinct genus of progymnosperm tree with fern-like leaves. A useful index fossil, this tree is found in strata dating from the Upper Devonian to Lower Carboniferous (), the oldest fossils being 385 million years old, and had global distribution. Until the 2007 discovery of Wattieza, many scientists considered Archaeopteris to be the earliest known tree. Bearing buds, reinforced branch joints, and branched trunks similar to today's woody plants, it is more reminiscent of modern seed-bearing trees than other spore-bearing taxa. It combines characteristics of woody trees and herbaceous ferns, and belongs to the progymnosperms, a group of extinct plants more closely related to seed plants than to ferns, but unlike seed plants, reproducing using spores like ferns. Taxonomy John William Dawson described the genus in 1871. The name derives from the ancient Greek (archaīos, "ancient"), and (ptéris, "fern"). Archaeopteris was originally classified as a fern, and it remained classified so for over 100 years. In 1911, Russian paleontologist Mikhail Dimitrievich Zalessky described a new type of petrified wood from the Donets Basin in modern Ukraine. He called the wood Callixylon, though he did not find any structures other than the trunk. The similarity to conifer wood was recognized. It was also noted that ferns of the genus Archaeopteris were often found associated with fossils of Callixylon. In the 1960s, paleontologist Charles B. Beck was able to demonstrate that the fossil wood known as Callixylon and the leaves known as Archaeopteris were actually part of the same plant. It was a plant with a mixture of characteristics not seen in any living plant, a link between true gymnosperms and ferns. The genus Archaeopteris is placed in the order Archaeopteridales and family Archaeopteridaceae. The name is similar to that of the first known feathered bird, Archaeopteryx, but in this case refers to the fern-like nature of the plant's fronds. Relationship to spermatophytes Archaeopteris is a member of a group of free-sporing woody plants called the progymnosperms that are interpreted as distant ancestors of the gymnosperms. Archaeopteris reproduced by releasing spores rather than by producing seeds, but some of the species, such as Archaeopteris halliana were heterosporous, producing two types of spores. This is thought to represent an early step in the evolution of vascular plants towards reproduction by seeds, which first appeared in the earliest, long-extinct gymnosperm group, the seed ferns (Pteridospermatophyta). The conifers or Pinophyta are one of four divisions of extant gymnosperms that arose from the seed ferns during the Carboniferous period. Description The trees of this genus typically grew to in height with leafy foliage reminiscent of some conifers. The large fern-like fronds were thickly set with fan-shaped leaflets or pinnae. The trunks of some species exceeded in diameter. The branches were borne in spiral arrangement, and a forked stipule was present at the base of each branch. Within a branch, leafy shoots were in opposite arrangement in a single plane. On fertile branches, some of the leaves were replaced by sporangia (spore capsules). Other modern adaptations Aside from its woody trunk, Archaeopteris possessed other modern adaptations to light interception and perhaps to seasonality as well. The large umbrella of fronds seems to have been quite optimized for light interception at the canopy level. In some species, the pinnules were shaped and oriented to avoid shading one another. There is evidence that whole fronds were shed together as single units, perhaps seasonally like modern deciduous foliage or like trees in the cypress family Cupressaceae. The plant had nodal zones that would have been important sites for the subsequent development of lateral roots and branches. Some branches were latent and adventitious, similar to those produced by living trees that eventually develop into roots. Before this time, shallow, rhizomatous roots had been the norm, but with Archaeopteris, deeper root systems were being developed that could support ever higher growth. Habitat Evidence indicates that Archaeopteris preferred wet soils, growing close to river systems and in floodplain woodlands. It would have formed a significant part of the canopy vegetation of early forests. Speaking of the first appearance of Archaeopteris on the world-scene, Stephen Scheckler, a professor of biology and geological sciences at Virginia Polytechnic Institute, says, "When [Archaeopteris] appears, it very quickly became the dominant tree all over the Earth. On all of the land areas that were habitable, they all had this tree". One species, Archaeopteris notosaria, has even been reported from within what was then the Antarctic Circle: leaves and fertile structures were identified from the Waterloo Farm lagerstätte in what is now South Africa. Scheckler believes that Archaeopteris had a major role in transforming its environment. "Its litter fed the streams and was a major factor in the evolution of freshwater fishes, whose numbers and varieties exploded in that time, and influenced the evolution of other marine ecosystems. It was the first plant to produce an extensive root system, so had a profound impact on soil chemistry. And once these ecosystem changes happened, they were changed for all time. It was a one-time thing." Looking roughly like a top-heavy Christmas tree, Archaeopteris may have played a part in the transformation of Earth's climate during the Devonian before becoming extinct within a short period of time at the beginning of the Carboniferous period.	Biology and health sciences	Pteridophytes	Plants
1495421	https://en.wikipedia.org/wiki/Li%20%28unit%29	Li (unit)	Li or ri (, lǐ, or , shìlǐ), also known as the Chinese mile, is a traditional Chinese unit of distance. The li has varied considerably over time but was usually about one third of an English mile and now has a standardized length of a half-kilometer (). This is then divided into 1,500 chi or "Chinese feet". The character 里 combines the characters for "field" (田, tián) and "earth" (土, tǔ), since it was considered to be about the length of a single village. As late as the 1940s, a "li" did not represent a fixed measure but could be longer or shorter depending on the effort required to cover the distance. There is also another li (Traditional: 釐, Simplified: 厘, lí) that indicates a unit of length of a chi, but it is used much less commonly. This li is used in the People's Republic of China as the equivalent of the centi- prefix in metric units, thus limi (厘米, límǐ) for centimeter. The tonal difference makes it distinguishable to speakers of Chinese, but unless specifically noted otherwise, any reference to li will always refer to the longer traditional unit and not to either the shorter unit or the kilometer. This traditional unit, in terms of historical usage and distance proportion, can be considered the East Asian counterpart to the Western league unit. However, in English league commonly means "3 miles." Changing values Like most traditional Chinese measurements, the li was reputed to have been established by the Yellow Emperor at the founding of Chinese civilization around 2600 BC and standardized by Yu the Great of the Xia dynasty six hundred years later. Although the value varied from state to state during the Spring and Autumn period and Warring States periods, historians give a general value to the li of 405 meters prior to the Qin dynasty imposition of its standard in the 3rd century BC. The basic Chinese traditional unit of distance was the chi. As its value changed over time, so did the lis. In addition, the number of chi per li was sometimes altered. To add further complexity, under the Qin dynasty, the li was set at 360 "paces" (, bù) but the number of chi per bu was subsequently changed from 6 to 5, shortening the li by . Thus, the Qin li of about 576 meters became (with other changes) the Han li, which was standardized at 415.8 meters. The basic units of measurement remained stable over the Qin and Han periods. A bronze imperial standard measure, dated AD 9, had been preserved at the Imperial Palace in Beijing and came to light in 1924. This has allowed very accurate conversions to modern measurements, which has provided a new and extremely useful additional tool in the identification of place names and routes. These measurements have been confirmed in many ways including the discovery of a number of rulers found at archaeological sites, and careful measurements of distances between known points. The Han li was calculated by Dubs to be 415.8 metres and all indications are that this is a precise and reliable determination. Under the Tang dynasty (AD 618–907), the li was approximately 323 meters. In the late Manchu or Qing dynasty, the number of chi was increased from 1,500 per li to 1,800. This had a value of 2115 feet or 644.6 meters. In addition, the Qing added a longer unit called the tu, which was equal to 150 li (96.7 km). These changes were undone by the Republic of China of Chiang Kai-shek, who adopted the metric system in 1928. The Republic of China (now also known as Taiwan) continues not to use the li at all but only the kilometer (Mandarin: , gōnglǐ, lit. "common li"). Under Mao Zedong, the People's Republic of China reinstituted the traditional units as a measure of anti-imperialism and cultural pride before officially adopting the metric system in 1984. A place was made within this for the traditional units, which were restandardized to metric values. A modern li is thus set at exactly half a kilometer (500 meters). However, unlike the jin which is still frequently preferred in daily use over the kilogram, the li is almost never used. Nonetheless, its appearance in many phrases and sayings means that "kilometer" must always be specified by saying gōnglǐ in full. Cultural use As one might expect for the equivalent of "mile", li appears in many Chinese sayings, locations, and proverbs as an indicator of great distances or the exotic: One Chinese name for the Great Wall is the "Ten-Thousand-Li Long Wall" (). As in Greek, the number "ten thousand" is used figuratively in Chinese to mean any "immeasurable" value and this title has never provided a literal distance of 10,000 li (). The actual length of the modern Great Wall is around 42,000 li (), over 4 times the name's proverbially "immeasurable" length. The Chinese proverb appearing in chapter 64 of the Tao Te Ching and commonly rendered as "A journey of a thousand miles begins with a single step" in fact refers to a thousand li: 千里之行，始於足下 (Qiānlǐzhīxíng, shǐyúzúxià). The greatest horses of Chinese history including Red Hare and Hualiu (驊騮) are all referred to as "thousand-li horses" (, qiānlǐmǎ), since they could supposedly travel a thousand li () in a single day. Li is sometimes used in location names, for example: Wulipu (Chinese: 五里铺镇), Hubei; Ankang Wulipu Airport (Chinese: 安康五里铺机场), Shaanxi. Sanlitun () is an area in Beijing. Ri in Japan and Korea The present day Korean ri (리, 里) and Japanese ri (里) are units of measurements that can be traced back to the Chinese li (里). Although the Chinese unit was unofficially used in Japan since the Zhou dynasty, the countries officially adopted the measurement used by the Tang dynasty (618–907 AD). The ri of an earlier era in Japan was thus true to Chinese length, corresponding to six chō ( 500–600 m), but later evolved to denote the distance that a person carrying a load would aim to cover on mountain roads in one hour. Thus, there had been various ri of 36, 40, and 48 chō. In the Edo period, the Tokugawa shogunate defined 1 ri as 36 chō, allowing other variants, and the Japanese government adopted this last definition in 1891. The Japanese ri was, at that time, fixed to the metric system, ≈ 3.93 kilometres or about 2.44 miles. Therefore, one must be careful about the correspondence between chō and ri. See Kujūkuri Beach (99-ri beach) for a case. In South Korea, the ri currently in use is a unit taken from the Han dynasty (206 BC–220 AD) li. It has a value of approximately 392.72 meters, or one tenth of the ri. The Aegukga, the national anthem of South Korea, and the Aegukka, the national anthem of North Korea, both mention 3,000 ri, which roughly corresponds to 1,200 km, the approximate longitudinal span of the Korean peninsula. In North Korea the Chollima Movement, a campaign aimed at improving labour productivity along the lines of the earlier Soviet Stakhanovite movement, gets its name from the word "chollima" which refers to a thousand-ri horse (chŏn + ri + ma in North Korean Romanization).	Physical sciences	East Asian	Basics and measurement
1495725	https://en.wikipedia.org/wiki/Picea%20rubens	Picea rubens	Picea rubens, commonly known as red spruce, is a species of spruce native to eastern North America, ranging from eastern Quebec and Nova Scotia, west to the Adirondack Mountains and south through New England along the Appalachians to western North Carolina and eastern Tennessee. This species is also known as yellow spruce, West Virginia spruce, eastern spruce, and he-balsam. Red spruce is the provincial tree of Nova Scotia. Description Red spruce is a perennial, shade-tolerant, late successional coniferous tree that under optimal conditions grows to tall with a trunk diameter of about , though exceptional specimens can reach tall and in diameter. It has a narrow conical crown. The leaves are needle-like, yellow-green, long, four-sided, curved, with a sharp point, and extend from all sides of the twig. The bark is gray-brown on the surface and red-brown on the inside, thin, and scaly. The wood is light, soft, has narrow rings, and has a slight red tinge. The cones are cylindrical, long, with a glossy red-brown color and stiff scales. The cones hang down from branches. Habitat Red spruce grows at a slow to moderate rate, lives for 250 to 450+ years, and is very shade-tolerant when young. It is often found in pure stands or forests mixed with eastern white pine, balsam fir, or black spruce. Along with Fraser fir, red spruce is one of two primary tree types in the southern Appalachian spruce-fir forest, a distinct ecosystem found only in the highest elevations of the southern Appalachian Mountains. Its habitat is moist but well-drained sandy loam, often at high altitudes. Red spruce can be easily damaged by windthrow and acid rain. Notable red spruce forests in West Virginia can be seen at Gaudineer Scenic Area, Canaan Valley, Roaring Plains West Wilderness, Dolly Sods Wilderness, and Spruce Mountain, all sites of former extensive red spruce forest. Related species It is closely related to black spruce, and hybrids between the two are frequent where their ranges meet. Genetic data suggests that the red spruce peripatrically speciated from the black spruce during glaciation in the Pleistocene. Uses Red spruce is used for Christmas trees and is an important wood used in making paper pulp. It is also an excellent tonewood and is used in many higher-end acoustic guitars and violins, as well as sound boards. The sap can be used to make spruce gum. Leafy red spruce twigs are boiled with sugar and flavoring to make spruce beer or spruce pudding. It can be used as construction lumber and is good for millwork and for crates. Damaging factors Like most trees, red spruce is subject to insect parasitism. Their insect enemy is the spruce budworm, although it is a bigger problem for white spruce and balsam fir. Other issues that have been damaging red spruce have included the increase in acid rain and current climate change. One of the consequences of acid rain deposition is the decrease of soil exchangeable calcium and increase of aluminum. This is because acid precipitation disrupts cation and nutrition cycling in forest ecosystems. Components of acid rain such as H+, NO3−, and SO42- limit the uptake of calcium by trees and can increase aluminum availability. Calcium concentration is important for red spruce for physiological processes such as dark respiration and cold tolerance, as well as disease resistance, signal transduction, membrane and cell wall synthesis and function, and regulation of stomata. Conversely, dissolved aluminum can be toxic or can interfere with root uptake of calcium and other nutrients. At the ecosystem and community levels, Calcium availability is associated with community composition, mature tree growth, and ecosystem productivity. One study testing the effects of added aluminum to soil, found that P. rubens mortality rate increased under these conditions. During the 1980s, increased acid deposition contributed to a loss of high-elevation red spruce trees caused by leached calcium and thus decreased freezing tolerance. Additionally, the structure of the spruce needle enhances the capture of water and particles, which has been shown to add to soil acidification, nutrient leaching, and forest decline. However, more recently, reductions in acid deposition have contributed to red spruce resurgence in some mountain areas in the northeastern United States. This increase in red spruce growth has been associated with an increase in rainfall pH, which reduces bulk acidic deposition. This suggests that policies aiming to reduce atmospheric pollution in this area have been effective, although other species sensitive to soil acidification, such as sugar maple, are still continuing to decline. Conservation The red spruce has low genetic diversity as well as a narrow ecological niche, meaning the tree is easily susceptible to changes within its environment. The Central Appalachian Spruce Restoration Initiative seeks to unite diverse partners with the goal of restoring historic red spruce ecosystems across the high-elevation landscapes of central Appalachians. The partners that make up this diverse group are Appalachian Mountain Joint Venture, Appalachian Regional Reforestation Initiative, Canaan Valley National Wildlife Refuge, Natural Resources Conservation Service, The Mountain Institute, The Nature Conservancy, Trout Unlimited, U.S. Forest Service Northern Research Station, U.S. Forest Service Monongahela National Forest, West Virginia Division of Natural Resources, West Virginia Division of Forestry, West Virginia Highlands Conservancy, West Virginia State Parks, and West Virginia University. Prior to the late 19th century, of red spruce were in West Virginia. In the late 19th and early 20th century, a vast amount of logging began in the state, and the number of red spruce dwindled to . Silviculture is being used to help restore the population of the lost red spruce. Significant efforts have been made to increase the growth of red spruce trees in western North Carolina, most notably by Molly Tartt on behalf of the Daughters of the American Revolution (DAR). Tartt embarked on a mission to find the lost red spruce forest that had been planted by the DAR as a memorial to the lives lost during the American Revolution. The forest, consisting of 50,000 trees was dedicated in 1940 and had until recently been forgotten until Tartt located and identified the forest near Devil's Courthouse.	Biology and health sciences	Pinaceae	Plants
1499520	https://en.wikipedia.org/wiki/Funnel%20cloud	Funnel cloud	A funnel cloud is a funnel-shaped cloud of condensed water droplets, associated with a rotating column of wind and extending from the base of a cloud (usually a cumulonimbus or towering cumulus cloud) but not reaching the ground or a water surface. A funnel cloud is usually visible as a cone-shaped or needle like protuberance from the main cloud base. Funnel clouds form most frequently in association with supercell thunderstorms, and are often, but not always, a visual precursor to tornadoes. Funnel clouds are visual phenomena, but these are not the vortex of wind itself. "Classic" funnel clouds If a funnel cloud touches the surface, the feature is considered a tornado, although ground level circulations begin before the visible condensation cloud appears. Most tornadoes begin as funnel clouds, but some funnel clouds do not make surface contact and these cannot be counted as tornadoes from the perspective of a naked eye observer, even as tornadic circulations of some intensity almost always are detectable when low-level radar observations are available. Also, tornadoes occur with some frequency without an associated condensation funnel. The term condensation funnel may refer to either a tornadic cloud or a funnel cloud aloft, but the term funnel cloud exclusively refers to a rotating condensation funnel not reaching the surface. If strong cyclonic winds are occurring at the surface and are connected to a cloud base, regardless of condensation, then the feature is a tornado. Funnel clouds result from the low air pressures found within tornadoes. The low air pressure causes air flowing towards the vortex to expand and cool. If the air is sufficiently moist and cools to the dew point, a funnel cloud is produced. The air pressure around the outer edge of the funnel cloud generally corresponds to the air pressure of the cloud base of the parent cloud. Debris swirls are usually evident prior to the condensation funnel reaching the surface. Some tornadoes may appear only as a debris swirl, with no obvious funnel cloud extending below the rotating cloud base at any time during the tornadic life cycle. The surface level vortex tends to strengthen over time following initial formation, making the debris swirls and the condensation more apparent. In cloud nomenclature, any funnel- or inverted-funnel-shaped cloud descending from cumulus or cumulonimbus clouds is technically described as an accessory feature called tuba. The terms tuba and funnel cloud are nearly but not exactly synonymous; a wall cloud, for example, is also a form of tuba. Funnel clouds associated with supercells usually form within and under wall clouds. Cold-air funnel clouds Cold-air funnel clouds (vortices) are generally much weaker than the vortices produced by supercells. Although cold-air funnels rarely make ground contact, surface level vortices sometimes become strong enough for condensation cloud to "touch down" briefly, becoming visible as weak tornadoes or waterspouts. Unlike the related phenomenon associated with severe thunderstorms, cold-air funnels are generally associated with partly cloudy skies in the wake of cold fronts, especially associated with certain low pressure systems, or in association with atmospheric boundaries such as lake and sea breezes or outflow boundaries. The larger scale weather conditions are characterized by particularly cold air aloft over relatively warmer low level air, leading to high lapse rates, and often by high environmental vorticity yet relatively meager vertical wind shear. The funnels develop where atmospheric instability and moisture are sufficient to support towering cumulus clouds but typically limited to no or to little precipitation. Cold-air funnels, although weak, may persist for several minutes, and areas of intermittently forming funnel clouds may occur for tens of minutes. Multiple such areas of activity may form within the same region during afternoon heating. Cold-air funnels appears to be a diurnal phenomenon, weakening and eventually dissipating with loss of solar heating. When precipitation does develop, the associated downdraft tends to cause rapid demise of the cold-air funnels. The mixing of cooler air in the lower troposphere with air flowing in a different direction in the middle troposphere causes the rotation on a horizontal axis, which, when deflected and tightened vertically by convective updrafts, forms a vertical rotation that can cause condensation to form a funnel cloud. Cold-air funnel clouds are a common sight along the Pacific Coast of the United States, particularly in the spring or fall. On July 29, 2013, a cold-core funnel cloud touched down as an EF0 tornado in Ottawa, Ontario, Canada, causing extensive damage in the form of downed trees on a golf course. No advance weather watches or warnings were issued by Environment Canada, and the tornado was spawned from one of the few non-severe storm clouds moving through the area. Other funnel clouds Other funnel clouds include shear or "high based" funnels, which are ephemeral, small, and weak funnels associated with small cumulus clouds, often even those rooted aloft above the boundary or surface layer, and in "fair weather" conditions. Small funnel clouds, such as some occurring within vicinity of mountains, occur by unknown processes. Shear funnels might also occur with weak transient circulations at the cloud base of thunderstorms. Mesoanticyclones, which accompany mesocyclones, often exhibit these funnel clouds. Brief funnels also are observed with some rear flank downdrafts (RFDs) (within inflow or outflow areas, and especially within inflow-outflow interchange areas as RFDs interact with mesocyclones or flanking line updrafts) and streamwise vorticity currents (SVCs) feeding into mesocyclones. Although not considered a separate kind of funnel cloud, some funnel clouds form with mesovortices associated with squall lines, which also can become tornadoes but are often not visible as funnel clouds or tornadoes because they usually occur within obscuring precipitation. Other "fair weather" funnel clouds include horseshoe clouds which are a very transient phenomena associated with extremely weak vortices.	Physical sciences	Clouds	Earth science
2145666	https://en.wikipedia.org/wiki/Toxodon	Toxodon	Toxodon (meaning "bow tooth" in reference to the curvature of the teeth) is an extinct genus of large ungulate native to South America from the Pliocene to the end of the Late Pleistocene. Toxodon is a member of Notoungulata, an order of extinct South American native ungulates distinct from the two living ungulate orders that had been indigenous to the continent for over 60 million years since the early Cenozoic, prior to the arrival of living ungulates into South America around 2.5 million years ago during the Great American Interchange. Toxodon is a member of the family Toxodontidae, which includes medium to large sized herbivores. Toxodon was one of the largest members of Toxodontidae and Notoungulata, with Toxodon platensis having an estimated body mass of . Remains of Toxodon were first collected by Charles Darwin during the voyage of the Beagle in 1832-33, and later scientifically named by Richard Owen in 1837. Both Darwin and Owen were puzzled by Toxodon's unusual anatomical features, including its long, ever-growing cheek teeth. Toxodon has been found across much of South America, excluding southern Patagonia, the Andes and northeastern-most region of the continent, inhabiting steppe, savanna and sometimes woodland habitats. Evidence suggests that Toxodon was ecologically plastic and able to adapt its diet to local conditions. While some authors have suggested that Toxodon was semiaquatic, isotopic analysis has supported a terrestrial lifestyle. Toxodon became extinct as part of the end-Pleistocene extinctions around 12,000 years ago, along with most large mammals across the Americas. The extinctions followed the arrival of humans to South America, who may have been a contributory factor in the extinctions. Several sites have been found suggesting human interaction with Toxodon. Taxonomy and evolution Charles Darwin, who was in South America as part of the second voyaging expedition of HMS Beagle, was one of the first to collect Toxodon fossils. In September–October 1832 and October 1833, Darwin collected several isolated teeth as well as a mandible from various localities in northern Argentina. On November 26, 1833, Darwin paid 18 pence (equivalent to £6.40 in 2018) for a T. platensis skull from a farmer in Uruguay. In his book covering the expedition, The Voyage of the Beagle. Darwin wrote, "November 26th – I set out on my return in a direct line for Montevideo. Having heard of some giant's bones at a neighbouring farm-house on the Sarandis, a small stream entering the Rio Negro, I rode there accompanied by my host, and purchased for the value of eighteen pence the head of the Toxodon." The skull had been propped up against a fence and been used as target practice for throwing stones by local children, who had knocked out its teeth. Since Darwin discovered that the fossils of similar mammals of South America were different from those in Europe, he invoked many debates about the evolution and natural selection of animals. In his own words, Darwin wrote down in his journal, Toxodon and its type species, T. platensis, were described in 1837 by Richard Owen based on remains collected by Darwin, in a paper titled "A description of the cranium of the Toxodon platensis, a gigantic extinct mammiferous species, referrible by its dentition to the Rodentia, but with affinities to the Pachydermata and the herbivorous Cetacea", reflecting the many unusual characteristics of its anatomy. Evolution Toxodon is a member of Notoungulata, a group of South American native ungulates that had been part of the fauna of South America since the Paleocene, over 60 million years ago, and had evolved in isolation in South America, prior to the arrival of living ungulates in South America around 2.5 million years ago as part of the Great American Interchange. Notoungulata represents the most diverse group of indigenous South American ungulates, with over 150 described genera in 13 different families. Notoungulates are morphologically diverse, including forms morphologically distant from Toxodon such as rodent and rabbit-like forms. Analysis of collagen sequences obtained from Toxodon as well as from Macrauchenia, a member of another indigenous South American ungulate order, Litopterna, found that notoungulates and litopterns were closely related to each other, and form a sister group to perissodactyls (which contains equids, rhinoceroses and tapirs) as part of the clade Panperissodactyla, making them true ungulates. This finding has been corroborated by an analysis of mitochondrial DNA extracted from a Macrauchenia fossil, which yielded a date of 66 million years ago for the time of the split from perissodactyls. Toxodon belongs to Toxodontidae, a large bodied group of notoungulates which first appeared in the Late Oligocene (Deseadan), ~28-23 million years ago, and underwent a great radiation during the Miocene epoch (~23-5.3 million years ago), when they reached their apex of diversity. The diversity of toxodontids, along with other notoungulates began to decline from around the Pliocene onwards, possibly as a result of climate change, as well as the arrival of competitors and predators from North America during the Great American Interchange following formation of the Isthmus of Panama. By the Late Pleistocene (Lujanian), the once great diversity of notoungulates had declined to only a few of species of toxodontids, with all other notoungulate families having become extinct. Cladogram of Toxodontidae, showing the position of Toxodon relative to other toxodontids, after Forasiepi et al., 2014: Species There has not been a recent taxonomic revision of the genus Toxodon, leaving the number of valid species uncertain. The species Toxodon chapalmalensis is known from the Pliocene (Montehermosan-Chapadmalalan) of Argentina, while Toxodon platensis, the type species, is known from the Pleistocene. The validity of other potential species like Toxodon darwini Burmeister, 1866, and Toxodon ensenadensis Ameghino, 1887 from the Early Pleistocene of Argentina is uncertain, and the species Toxodon gezi C. Ameghino, 1917 and Toxodon aguirrei Ameghino, 1917 have been considered junior synonyms of Toxodon platensis by recent authors. Some recent authors have argued that Toxodon gracilis Gervais and Ameghino, 1880, should be recognised as a distinct species from the Pleistocene of the Pampas significantly smaller than T. platensis, with these authors suggesting that T. platensis and T. gracilis represent the only valid species of Toxodon in the Pleistocene of the Pampas region. Other authors have argued that all Pleistocene Toxodon species should be considered synonymous with T. platensis. Description The bodyform of Toxodon and other toxodontids have been compared to those of hippopotamuses and rhinoceroses. Toxodon platensis is one of the largest known toxodontids and notoungulates, with an estimated body mass of approximately , and a body length of around . The skull of Toxodon is proportionally large, and triangular in shape when viewed from above. All of the teeth in the jaws are high-crowned (hypsodont). Like other toxodontids, the upper and lower first incisors (I1 and i1) are large and protrude, with the second upper incisors (I2) and lower third incisors (i3) being modified into evergrowing tusks. The upper incisors display an arched shape, while the lower incisors project horizontally forwards at the front of the lower jaw. The wide front of the lower jaw with the horizontally-arranged incisors has been described as "spade-like". There is a gap (diastema) between the incisors and the cheek teeth. Like other derived toxodontids, Toxodon had long, ever-growing (hypselodont) cheek (premolar and molar) teeth, with the name Toxodon deriving from the curved shape of the upper molars, which are bowed inwards towards the midline of the skull to fit in the upper jaw. Evergrowing cheek teeth are unknown in any living ungulates, but are present in some other mammal groups like wombats and rodents. The surface of the cheek teeth is primarily composed of dentine. The thoracic vertebrae of Toxodon have elongate neural spines, which likely anchored muscles which supported the large head. The legs of Toxodon are relatively short, with their bones being robust. The hindlimb is considerably longer than the forelimb, resulting in the back being elevated and the shoulder, neck and head being relatively low. The ulna has a strongly backwardly projecting olecranon process similar to that of rhinos, suggesting that the front leg was held extended when standing. The (distal) part of the femur closest to the foot shows a pronounced medial trochlear ridge, which likely served along with the patella (kneecap) to allow the knees to be locked when standing akin to the stay apparatus of living horses as an energy saving mechanism. There are three functional digits on each foot, which are tipped with hoof-like phalanges. Distribution Toxodon had a widespread distribution in South America east of the Andes, ranging from northern Argentina and Bolivia to the western Amazon on the Peru-Brazil border, to Northeast Brazil. Although some authors suggest that the distribution of Toxodon extended into Venezuela, other authors suggest that the related Mixotoxodon (which ranged as far north as the southern United States) was the only toxodontid present in the region during the Pleistocene. Palaeobiology Although some authors have suggested that Toxodon was semiaquatic based on the similarity of some aspects of its anatomy to hippopotamuses, this has been disputed by other authors, and analysis of oxygen isotope ratios (which differ between terrestrial and aquatic animals) suggests a terrestrial lifestyle for Toxodon. As such, it has been suggested that Toxodon was probably more ecologically comparable to rhinoceroses. Toxodon is suggested to have been capable of moving at considerable speed. Toxodon is believed to have been ecologically plastic and have had a wide niche breadth, with its diet varying according to local conditions, with an almost totally C3 browsing diet in the Amazon rainforest, mixed feeding C3 in Bahia and the Pampas to almost completely C4 dominated grazing diet in the Chaco. Within the Brazilian Intertropical Region, local climate had little impact on the diet of T. platensis. Although Toxodon is thought to have inhabited open landscapes like steppe and savannah, in some areas like the southwestern Brazilian Amazon, it is suggested to have inhabited woodland. Like living animals of similar size, it has been suggested that Toxodon probably only gave birth to a single offspring at a time. T. platensis bones have been found displaying signs of disease like osteomyelitis and spondyloarthropathies. The teeth of Toxodon often display enamel hypoplasia (loss of tooth enamel) in the form of grooves and pits, which is likely due to their evergrowing nature and/or environmental stresses. Tracks probably attributable to Toxodon have been reported from eastern Pernambuco in Northeast Brazil. Isotopic analysis suggests that Toxodon may have been predated upon by the large sabertooth cat Smilodon populator, the apex predator of South American ecosystems during much of the Pleistocene. Extinction Toxodon became extinct at the end of the Late Pleistocene around 12,000 years as part of the end-Pleistocene extinction event alongside almost all other large animals in South America. Previous mid-Holocene dates are now thought to be in error. These extinctions followed the first arrival of humans in the Americas, and it has been suggested human hunting may have been a casual factor in the extinctions. Several sites record apparent interactions between Toxodon and humans. Remains of Toxodon from the Arroyo Seco 2 site in the Pampas are associated with unambiguously butchered megafauna, but it is unclear if the Toxodon itself was actually butchered or the remains were naturally transported to the site. At the Paso Otero 5 site in the Pampas of northeast Argentina, burned bones of Toxodon alongside those of numerous other extinct megafauna species are associated with Fishtail points (a type of knapped stone spear point common across South America at the end of the Pleistocene, suggested to be used to hunt large mammals). The bones of the megafauna were probably deliberately burned as fuel. No cut marks are visible on the vast majority of bones at the site (with only one bone of a llama possibly displaying any butchery marks), which may be due to the burning degrading the bones. Various remains of Toxodon platensis in the collection of the Museum national d'Histoire naturelle collected from the Pampas region in the 19th century including a femur, an iliac fragment, a tibia, as well as a mandible (the latter of which has been radiocarbon dated to around 13,000 years ago), have been found to display cut marks indicative of butchery.	Biology and health sciences	Mammals: General	Animals
2146034	https://en.wikipedia.org/wiki/CRISPR	CRISPR	CRISPR () (an acronym for clustered regularly interspaced short palindromic repeats) is a family of DNA sequences found in the genomes of prokaryotic organisms such as bacteria and archaea. Each sequence within an individual prokaryotic cell is derived from a DNA fragment of a bacteriophage that had previously infected the prokaryote or one of its ancestors. These sequences are used to detect and destroy DNA from similar bacteriophages during subsequent infections. Hence these sequences play a key role in the antiviral (i.e. anti-phage) defense system of prokaryotes and provide a form of heritable, acquired immunity. CRISPR is found in approximately 50% of sequenced bacterial genomes and nearly 90% of sequenced archaea. Cas9 (or "CRISPR-associated protein 9") is an enzyme that uses CRISPR sequences as a guide to recognize and open up specific strands of DNA that are complementary to the CRISPR sequence. Cas9 enzymes together with CRISPR sequences form the basis of a technology known as CRISPR-Cas9 that can be used to edit genes within living organisms. This editing process has a wide variety of applications including basic biological research, development of biotechnological products, and treatment of diseases. The development of the CRISPR-Cas9 genome editing technique was recognized by the Nobel Prize in Chemistry in 2020 awarded to Emmanuelle Charpentier and Jennifer Doudna. History Repeated sequences The discovery of clustered DNA repeats took place independently in three parts of the world. The first description of what would later be called CRISPR is from Osaka University researcher Yoshizumi Ishino and his colleagues in 1987. They accidentally cloned part of a CRISPR sequence together with the "iap" gene (isozyme conversion of alkaline phosphatase) from their target genome, that of Escherichia coli. The organization of the repeats was unusual. Repeated sequences are typically arranged consecutively, without interspersing different sequences. They did not know the function of the interrupted clustered repeats. In 1993, researchers of Mycobacterium tuberculosis in the Netherlands published two articles about a cluster of interrupted direct repeats (DR) in that bacterium. They recognized the diversity of the sequences that intervened in the direct repeats among different strains of M. tuberculosis and used this property to design a typing method called spoligotyping, still in use today. Francisco Mojica at the University of Alicante in Spain studied the function of repeats in the archaeal species Haloferax and Haloarcula. Mojica's supervisor surmised that the clustered repeats had a role in correctly segregating replicated DNA into daughter cells during cell division, because plasmids and chromosomes with identical repeat arrays could not coexist in Haloferax volcanii. Transcription of the interrupted repeats was also noted for the first time; this was the first full characterization of CRISPR. By 2000, Mojica and his students, after an automated search of published genomes, identified interrupted repeats in 20 species of microbes as belonging to the same family. Because those sequences were interspaced, Mojica initially called these sequences "short regularly spaced repeats" (SRSR). In 2001, Mojica and Ruud Jansen, who were searching for an additional interrupted repeats, proposed the acronym CRISPR (Clustered Regularly Interspaced Short Palindromic Repeats) to unify the numerous acronyms used to describe these sequences. In 2002, Tang, et al. showed evidence that CRISPR repeat regions from the genome of Archaeoglobus fulgidus were transcribed into long RNA molecules subsequently processed into unit-length small RNAs, plus some longer forms of 2, 3, or more spacer-repeat units. In 2005, yogurt researcher Rodolphe Barrangou discovered that Streptococcus thermophilus, after iterative phage infection challenges, develops increased phage resistance due to the incorporation of additional CRISPR spacer sequences. Barrangou's employer, the Danish food company Danisco, then developed phage-resistant S. thermophilus strains for yogurt production. Danisco was later bought by DuPont, which owns about 50 percent of the global dairy culture market, and the technology spread widely. CRISPR-associated systems A major advance in understanding CRISPR came with Jansen's observation that the prokaryote repeat cluster was accompanied by four homologous genes that make up CRISPR-associated systems, cas 1–4. The Cas proteins showed helicase and nuclease motifs, suggesting a role in the dynamic structure of the CRISPR loci. In this publication, the acronym CRISPR was used as the universal name of this pattern, but its function remained enigmatic. In 2005, three independent research groups showed that some CRISPR spacers are derived from phage DNA and extrachromosomal DNA such as plasmids. In effect, the spacers are fragments of DNA gathered from viruses that previously attacked the cell. The source of the spacers was a sign that the CRISPR-cas system could have a role in adaptive immunity in bacteria. All three studies proposing this idea were initially rejected by high-profile journals, but eventually appeared in other journals. The first publication proposing a role of CRISPR-Cas in microbial immunity, by Mojica and collaborators at the University of Alicante, predicted a role for the RNA transcript of spacers on target recognition in a mechanism that could be analogous to the RNA interference system used by eukaryotic cells. Koonin and colleagues extended this RNA interference hypothesis by proposing mechanisms of action for the different CRISPR-Cas subtypes according to the predicted function of their proteins. Experimental work by several groups revealed the basic mechanisms of CRISPR-Cas immunity. In 2007, the first experimental evidence that CRISPR was an adaptive immune system was published. A CRISPR region in Streptococcus thermophilus acquired spacers from the DNA of an infecting bacteriophage. The researchers manipulated the resistance of S. thermophilus to different types of phages by adding and deleting spacers whose sequence matched those found in the tested phages. In 2008, Brouns and Van der Oost identified a complex of Cas proteins called Cascade, that in E. coli cut the CRISPR RNA precursor within the repeats into mature spacer-containing RNA molecules called CRISPR RNA (crRNA), which remained bound to the protein complex. Moreover, it was found that Cascade, crRNA and a helicase/nuclease (Cas3) were required to provide a bacterial host with immunity against infection by a DNA virus. By designing an anti-virus CRISPR, they demonstrated that two orientations of the crRNA (sense/antisense) provided immunity, indicating that the crRNA guides were targeting dsDNA. That year Marraffini and Sontheimer confirmed that a CRISPR sequence of S. epidermidis targeted DNA and not RNA to prevent conjugation. This finding was at odds with the proposed RNA-interference-like mechanism of CRISPR-Cas immunity, although a CRISPR-Cas system that targets foreign RNA was later found in Pyrococcus furiosus. A 2010 study showed that CRISPR-Cas cuts strands of both phage and plasmid DNA in S. thermophilus. Cas9 A simpler CRISPR system from Streptococcus pyogenes relies on the protein Cas9. The Cas9 endonuclease is a four-component system that includes two small molecules: crRNA and trans-activating CRISPR RNA (tracrRNA). In 2012, Jennifer Doudna and Emmanuelle Charpentier re-engineered the Cas9 endonuclease into a more manageable two-component system by fusing the two RNA molecules into a "single-guide RNA" that, when combined with Cas9, could find and cut the DNA target specified by the guide RNA. This contribution was so significant that it was recognized by the Nobel Prize in Chemistry in 2020. By manipulating the nucleotide sequence of the guide RNA, the artificial Cas9 system could be programmed to target any DNA sequence for separation. Another collaboration comprising Virginijus Šikšnys, Gasiūnas, Barrangou, and Horvath showed that Cas9 from the S. thermophilus CRISPR system can also be reprogrammed to target a site of their choosing by changing the sequence of its crRNA. These advances fueled efforts to edit genomes with the modified CRISPR-Cas9 system. Groups led by Feng Zhang and George Church simultaneously published descriptions of genome editing in human cell cultures using CRISPR-Cas9 for the first time. It has since been used in a wide range of organisms, including baker's yeast (Saccharomyces cerevisiae), the opportunistic pathogen Candida albicans, zebrafish (Danio rerio), fruit flies (Drosophila melanogaster), ants (Harpegnathos saltator and Ooceraea biroi), mosquitoes (Aedes aegypti), nematodes (Caenorhabditis elegans), plants, mice (Mus musculus domesticus), monkeys and human embryos. CRISPR has been modified to make programmable transcription factors that allows activation or silencing of targeted genes. The CRISPR-Cas9 system has been shown to make effective gene edits in Human tripronuclear zygotes, as first described in a 2015 paper by Chinese scientists P. Liang and Y. Xu. The system made a successful cleavage of mutant Beta-Hemoglobin (HBB) in 28 out of 54 embryos. Four out of the 28 embryos were successfully recombined using a donor template. The scientists showed that during DNA recombination of the cleaved strand, the homologous endogenous sequence HBD competes with the exogenous donor template. DNA repair in human embryos is much more complicated and particular than in derived stem cells. Cas12a In 2015, the nuclease Cas12a (formerly called ) was characterized in the CRISPR-Cpf1 system of the bacterium Francisella novicida. Its original name, from a TIGRFAMs protein family definition built in 2012, reflects the prevalence of its CRISPR-Cas subtype in the Prevotella and Francisella lineages. Cas12a showed several key differences from Cas9 including: causing a 'staggered' cut in double stranded DNA as opposed to the 'blunt' cut produced by Cas9, relying on a 'T rich' PAM (providing alternative targeting sites to Cas9), and requiring only a CRISPR RNA (crRNA) for successful targeting. By contrast, Cas9 requires both crRNA and a trans-activating crRNA (tracrRNA). These differences may give Cas12a some advantages over Cas9. For example, Cas12a's small crRNAs are ideal for multiplexed genome editing, as more of them can be packaged in one vector than can Cas9's sgRNAs. The sticky 5′ overhangs left by Cas12a can also be used for DNA assembly that is much more target-specific than traditional restriction enzyme cloning. Finally, Cas12a cleaves DNA 18–23 base pairs downstream from the PAM site. This means there is no disruption to the recognition sequence after repair, and so Cas12a enables multiple rounds of DNA cleavage. By contrast, since Cas9 cuts only 3 base pairs upstream of the PAM site, the NHEJ pathway results in indel mutations that destroy the recognition sequence, thereby preventing further rounds of cutting. In theory, repeated rounds of DNA cleavage should cause an increased opportunity for the desired genomic editing to occur. A distinctive feature of Cas12a, as compared to Cas9, is that after cutting its target, Cas12a remains bound to the target and then cleaves other ssDNA molecules non-discriminately. This property is called "collateral cleavage" or "trans-cleavage" activity and has been exploited for the development of various diagnostic technologies. Cas13 In 2016, the nuclease (formerly known as ) from the bacterium Leptotrichia shahii was characterized. Cas13 is an RNA-guided RNA endonuclease, which means that it does not cleave DNA, but only single-stranded RNA. Cas13 is guided by its crRNA to a ssRNA target and binds and cleaves the target. Similar to Cas12a, the Cas13 remains bound to the target and then cleaves other ssRNA molecules non-discriminately. This collateral cleavage property has been exploited for the development of various diagnostic technologies. In 2021, Dr. Hui Yang characterized novel miniature Cas13 protein (mCas13) variants, Cas13X and Cas13Y. Using a small portion of N gene sequence from SARS-CoV-2 as a target in characterization of mCas13, revealed the sensitivity and specificity of mCas13 coupled with RT-LAMP for detection of SARS-CoV-2 in both synthetic and clinical samples over other available standard tests like RT-qPCR (1 copy/μL). Locus structure Repeats and spacers The CRISPR array is made up of an AT-rich leader sequence followed by short repeats that are separated by unique spacers. CRISPR repeats typically range in size from 28 to 37 base pairs (bps), though there can be as few as 23 bp and as many as 55 bp. Some show dyad symmetry, implying the formation of a secondary structure such as a stem-loop ('hairpin') in the RNA, while others are designed to be unstructured. The size of spacers in different CRISPR arrays is typically 32 to 38 bp (range 21 to 72 bp). New spacers can appear rapidly as part of the immune response to phage infection. There are usually fewer than 50 units of the repeat-spacer sequence in a CRISPR array. CRISPR RNA structures Cas genes and CRISPR subtypes Small clusters of cas genes are often located next to CRISPR repeat-spacer arrays. Collectively the 93 cas genes are grouped into 35 families based on sequence similarity of the encoded proteins. 11 of the 35 families form the cas core, which includes the protein families Cas1 through Cas9. A complete CRISPR-Cas locus has at least one gene belonging to the cas core. CRISPR-Cas systems fall into two classes. Class 1 systems use a complex of multiple Cas proteins to degrade foreign nucleic acids. Class 2 systems use a single large Cas protein for the same purpose. Class 1 is divided into types I, III, and IV; class 2 is divided into types II, V, and VI. The 6 system types are divided into 33 subtypes. Each type and most subtypes are characterized by a "signature gene" found almost exclusively in the category. Classification is also based on the complement of cas genes that are present. Most CRISPR-Cas systems have a Cas1 protein. The phylogeny of Cas1 proteins generally agrees with the classification system, but exceptions exist due to module shuffling. Many organisms contain multiple CRISPR-Cas systems suggesting that they are compatible and may share components. The sporadic distribution of the CRISPR-Cas subtypes suggests that the CRISPR-Cas system is subject to horizontal gene transfer during microbial evolution. Mechanism CRISPR-Cas immunity is a natural process of bacteria and archaea. CRISPR-Cas prevents bacteriophage infection, conjugation and natural transformation by degrading foreign nucleic acids that enter the cell. Spacer acquisition When a microbe is invaded by a bacteriophage, the first stage of the immune response is to capture phage DNA and insert it into a CRISPR locus in the form of a spacer. Cas1 and Cas2 are found in both types of CRISPR-Cas immune systems, which indicates that they are involved in spacer acquisition. Mutation studies confirmed this hypothesis, showing that removal of Cas1 or Cas2 stopped spacer acquisition, without affecting CRISPR immune response. Multiple Cas1 proteins have been characterised and their structures resolved. Cas1 proteins have diverse amino acid sequences. However, their crystal structures are similar and all purified Cas1 proteins are metal-dependent nucleases/integrases that bind to DNA in a sequence-independent manner. Representative Cas2 proteins have been characterised and possess either (single strand) ssRNA- or (double strand) dsDNA- specific endoribonuclease activity. In the I-E system of E. coli Cas1 and Cas2 form a complex where a Cas2 dimer bridges two Cas1 dimers. In this complex Cas2 performs a non-enzymatic scaffolding role, binding double-stranded fragments of invading DNA, while Cas1 binds the single-stranded flanks of the DNA and catalyses their integration into CRISPR arrays. New spacers are usually added at the beginning of the CRISPR next to the leader sequence creating a chronological record of viral infections. In E. coli a histone like protein called integration host factor (IHF), which binds to the leader sequence, is responsible for the accuracy of this integration. IHF also enhances integration efficiency in the type I-F system of Pectobacterium atrosepticum. but in other systems, different host factors may be required Protospacer adjacent motifs (PAM) Bioinformatic analysis of regions of phage genomes that were excised as spacers (termed protospacers) revealed that they were not randomly selected but instead were found adjacent to short (3–5 bp) DNA sequences termed protospacer adjacent motifs (PAM). Analysis of CRISPR-Cas systems showed PAMs to be important for type I and type II, but not type III systems during acquisition. In type I and type II systems, protospacers are excised at positions adjacent to a PAM sequence, with the other end of the spacer cut using a ruler mechanism, thus maintaining the regularity of the spacer size in the CRISPR array. The conservation of the PAM sequence differs between CRISPR-Cas systems and appears to be evolutionarily linked to Cas1 and the leader sequence. New spacers are added to a CRISPR array in a directional manner, occurring preferentially, but not exclusively, adjacent to the leader sequence. Analysis of the type I-E system from E. coli demonstrated that the first direct repeat adjacent to the leader sequence is copied, with the newly acquired spacer inserted between the first and second direct repeats. The PAM sequence appears to be important during spacer insertion in type I-E systems. That sequence contains a strongly conserved final nucleotide (nt) adjacent to the first nt of the protospacer. This nt becomes the final base in the first direct repeat. This suggests that the spacer acquisition machinery generates single-stranded overhangs in the second-to-last position of the direct repeat and in the PAM during spacer insertion. However, not all CRISPR-Cas systems appear to share this mechanism as PAMs in other organisms do not show the same level of conservation in the final position. It is likely that in those systems, a blunt end is generated at the very end of the direct repeat and the protospacer during acquisition. Insertion variants Analysis of Sulfolobus solfataricus CRISPRs revealed further complexities to the canonical model of spacer insertion, as one of its six CRISPR loci inserted new spacers randomly throughout its CRISPR array, as opposed to inserting closest to the leader sequence. Multiple CRISPRs contain many spacers to the same phage. The mechanism that causes this phenomenon was discovered in the type I-E system of E. coli. A significant enhancement in spacer acquisition was detected where spacers already target the phage, even mismatches to the protospacer. This 'priming' requires the Cas proteins involved in both acquisition and interference to interact with each other. Newly acquired spacers that result from the priming mechanism are always found on the same strand as the priming spacer. This observation led to the hypothesis that the acquisition machinery slides along the foreign DNA after priming to find a new protospacer. Biogenesis CRISPR-RNA (crRNA), which later guides the Cas nuclease to the target during the interference step, must be generated from the CRISPR sequence. The crRNA is initially transcribed as part of a single long transcript encompassing much of the CRISPR array. This transcript is then cleaved by Cas proteins to form crRNAs. The mechanism to produce crRNAs differs among CRISPR-Cas systems. In type I-E and type I-F systems, the proteins Cas6e and Cas6f respectively, recognise stem-loops created by the pairing of identical repeats that flank the crRNA. These Cas proteins cleave the longer transcript at the edge of the paired region, leaving a single crRNA along with a small remnant of the paired repeat region. Type III systems also use Cas6, however, their repeats do not produce stem-loops. Cleavage instead occurs by the longer transcript wrapping around the Cas6 to allow cleavage just upstream of the repeat sequence. Type II systems lack the Cas6 gene and instead utilize RNaseIII for cleavage. Functional type II systems encode an extra small RNA that is complementary to the repeat sequence, known as a trans-activating crRNA (tracrRNA). Transcription of the tracrRNA and the primary CRISPR transcript results in base pairing and the formation of dsRNA at the repeat sequence, which is subsequently targeted by RNaseIII to produce crRNAs. Unlike the other two systems, the crRNA does not contain the full spacer, which is instead truncated at one end. CrRNAs associate with Cas proteins to form ribonucleotide complexes that recognize foreign nucleic acids. CrRNAs show no preference between the coding and non-coding strands, which is indicative of an RNA-guided DNA-targeting system. The type I-E complex (commonly referred to as Cascade) requires five Cas proteins bound to a single crRNA. Interference During the interference stage in type I systems, the PAM sequence is recognized on the crRNA-complementary strand and is required along with crRNA annealing. In type I systems correct base pairing between the crRNA and the protospacer signals a conformational change in Cascade that recruits Cas3 for DNA degradation. Type II systems rely on a single multifunctional protein, Cas9, for the interference step. Cas9 requires both the crRNA and the tracrRNA to function and cleave DNA using its dual HNH and RuvC/RNaseH-like endonuclease domains. Basepairing between the PAM and the phage genome is required in type II systems. However, the PAM is recognized on the same strand as the crRNA (the opposite strand to type I systems). Type III systems, like type I require six or seven Cas proteins binding to crRNAs. The type III systems analysed from S. solfataricus and P. furiosus both target the mRNA of phages rather than phage DNA genome, which may make these systems uniquely capable of targeting RNA-based phage genomes. Type III systems were also found to target DNA in addition to RNA using a different Cas protein in the complex, Cas10. The DNA cleavage was shown to be transcription dependent. The mechanism for distinguishing self from foreign DNA during interference is built into the crRNAs and is therefore likely common to all three systems. Throughout the distinctive maturation process of each major type, all crRNAs contain a spacer sequence and some portion of the repeat at one or both ends. It is the partial repeat sequence that prevents the CRISPR-Cas system from targeting the chromosome as base pairing beyond the spacer sequence signals self and prevents DNA cleavage. RNA-guided CRISPR enzymes are classified as type V restriction enzymes. Evolution The cas genes in the adaptor and effector modules of the CRISPR-Cas system are believed to have evolved from two different ancestral modules. A transposon-like element called casposon encoding the Cas1-like integrase and potentially other components of the adaptation module was inserted next to the ancestral effector module, which likely functioned as an independent innate immune system. The highly conserved cas1 and cas2 genes of the adaptor module evolved from the ancestral module while a variety of class 1 effector cas genes evolved from the ancestral effector module. The evolution of these various class 1 effector module cas genes was guided by various mechanisms, such as duplication events. On the other hand, each type of class 2 effector module arose from subsequent independent insertions of mobile genetic elements. These mobile genetic elements took the place of the multiple gene effector modules to create single gene effector modules that produce large proteins which perform all the necessary tasks of the effector module. The spacer regions of CRISPR-Cas systems are taken directly from foreign mobile genetic elements and thus their long-term evolution is hard to trace. The non-random evolution of these spacer regions has been found to be highly dependent on the environment and the particular foreign mobile genetic elements it contains. CRISPR-Cas can immunize bacteria against certain phages and thus halt transmission. For this reason, Koonin described CRISPR-Cas as a Lamarckian inheritance mechanism. However, this was disputed by a critic who noted, "We should remember [Lamarck] for the good he contributed to science, not for things that resemble his theory only superficially. Indeed, thinking of CRISPR and other phenomena as Lamarckian only obscures the simple and elegant way evolution really works". But as more recent studies have been conducted, it has become apparent that the acquired spacer regions of CRISPR-Cas systems are indeed a form of Lamarckian evolution because they are genetic mutations that are acquired and then passed on. On the other hand, the evolution of the Cas gene machinery that facilitates the system evolves through classic Darwinian evolution. Coevolution Analysis of CRISPR sequences revealed coevolution of host and viral genomes. The basic model of CRISPR evolution is newly incorporated spacers driving phages to mutate their genomes to avoid the bacterial immune response, creating diversity in both the phage and host populations. To resist a phage infection, the sequence of the CRISPR spacer must correspond perfectly to the sequence of the target phage gene. Phages can continue to infect their hosts' given point mutations in the spacer. Similar stringency is required in PAM or the bacterial strain remains phage sensitive. Rates A study of 124 S. thermophilus strains showed that 26% of all spacers were unique and that different CRISPR loci showed different rates of spacer acquisition. Some CRISPR loci evolve more rapidly than others, which allowed the strains' phylogenetic relationships to be determined. A comparative genomic analysis showed that E. coli and S. enterica evolve much more slowly than S. thermophilus. The latter's strains that diverged 250,000 years ago still contained the same spacer complement. Metagenomic analysis of two acid-mine-drainage biofilms showed that one of the analyzed CRISPRs contained extensive deletions and spacer additions versus the other biofilm, suggesting a higher phage activity/prevalence in one community than the other. In the oral cavity, a temporal study determined that 7–22% of spacers were shared over 17 months within an individual while less than 2% were shared across individuals. From the same environment, a single strain was tracked using PCR primers specific to its CRISPR system. Broad-level results of spacer presence/absence showed significant diversity. However, this CRISPR added three spacers over 17 months, suggesting that even in an environment with significant CRISPR diversity some loci evolve slowly. CRISPRs were analysed from the metagenomes produced for the Human Microbiome Project. Although most were body-site specific, some within a body site are widely shared among individuals. One of these loci originated from streptococcal species and contained ≈15,000 spacers, 50% of which were unique. Similar to the targeted studies of the oral cavity, some showed little evolution over time. CRISPR evolution was studied in chemostats using S. thermophilus to directly examine spacer acquisition rates. In one week, S. thermophilus strains acquired up to three spacers when challenged with a single phage. During the same interval, the phage developed single-nucleotide polymorphisms that became fixed in the population, suggesting that targeting had prevented phage replication absent these mutations. Another S. thermophilus experiment showed that phages can infect and replicate in hosts that have only one targeting spacer. Yet another showed that sensitive hosts can exist in environments with high-phage titres. The chemostat and observational studies suggest many nuances to CRISPR and phage (co)evolution. Identification CRISPRs are widely distributed among bacteria and archaea and show some sequence similarities. Their most notable characteristic is their repeating spacers and direct repeats. This characteristic makes CRISPRs easily identifiable in long sequences of DNA, since the number of repeats decreases the likelihood of a false positive match. Analysis of CRISPRs in metagenomic data is more challenging, as CRISPR loci do not typically assemble, due to their repetitive nature or through strain variation, which confuses assembly algorithms. Where many reference genomes are available, polymerase chain reaction (PCR) can be used to amplify CRISPR arrays and analyse spacer content. However, this approach yields information only for specifically targeted CRISPRs and for organisms with sufficient representation in public databases to design reliable polymerase PCR primers. Degenerate repeat-specific primers can be used to amplify CRISPR spacers directly from environmental samples; amplicons containing two or three spacers can be then computationally assembled to reconstruct long CRISPR arrays. The alternative is to extract and reconstruct CRISPR arrays from shotgun metagenomic data. This is computationally more difficult, particularly with second generation sequencing technologies (e.g. 454, Illumina), as the short read lengths prevent more than two or three repeat units appearing in a single read. CRISPR identification in raw reads has been achieved using purely de novo identification or by using direct repeat sequences in partially assembled CRISPR arrays from contigs (overlapping DNA segments that together represent a consensus region of DNA) and direct repeat sequences from published genomes as a hook for identifying direct repeats in individual reads. Use by phages Another way for bacteria to defend against phage infection is by having chromosomal islands. A subtype of chromosomal islands called phage-inducible chromosomal island (PICI) is excised from a bacterial chromosome upon phage infection and can inhibit phage replication. PICIs are induced, excised, replicated, and finally packaged into small capsids by certain staphylococcal temperate phages. PICIs use several mechanisms to block phage reproduction. In the first mechanism, PICI-encoded Ppi differentially blocks phage maturation by binding or interacting specifically with phage TerS, hence blocking phage TerS/TerL complex formation responsible for phage DNA packaging. In the second mechanism PICI CpmAB redirects the phage capsid morphogenetic protein to make 95% of SaPI-sized capsid and phage DNA can package only 1/3rd of their genome in these small capsids and hence become nonviable phage. The third mechanism involves two proteins, PtiA and PtiB, that target the LtrC, which is responsible for the production of virion and lysis proteins. This interference mechanism is modulated by a modulatory protein, PtiM, binds to one of the interference-mediating proteins, PtiA, and hence achieves the required level of interference. One study showed that lytic ICP1 phage, which specifically targets Vibrio cholerae serogroup O1, has acquired a CRISPR-Cas system that targets a V. cholera PICI-like element. The system has 2 CRISPR loci and 9 Cas genes. It seems to be homologous to the I-F system found in Yersinia pestis. Moreover, like the bacterial CRISPR-Cas system, ICP1 CRISPR-Cas can acquire new sequences, which allows phage and host to co-evolve. Certain archaeal viruses were shown to carry mini-CRISPR arrays containing one or two spacers. It has been shown that spacers within the virus-borne CRISPR arrays target other viruses and plasmids, suggesting that mini-CRISPR arrays represent a mechanism of heterotypic superinfection exclusion and participate in interviral conflicts. Applications CRISPR gene editing is a revolutionary technology that allows for precise, targeted modifications to the DNA of living organisms. Developed from a natural defense mechanism found in bacteria, CRISPR-Cas9 is the most commonly used system, that allows "cutting" of DNA at specific locations and either delete, modify, or insert genetic material. This technology has transformed fields such as genetics, medicine, and agriculture, offering potential treatments for genetic disorders, advancements in crop engineering, and research into the fundamental workings of life. However, its ethical implications and potential unintended consequences have sparked significant debate.	Technology	Biotechnology	null
2146163	https://en.wikipedia.org/wiki/SN%201572	SN 1572	SN 1572 (Tycho's Star, Tycho's Nova, Tycho's Supernova), or B Cassiopeiae (B Cas), was a supernova of Type Ia in the constellation Cassiopeia, one of eight supernovae visible to the naked eye in historical records. It appeared in early November 1572 and was independently discovered by many individuals. Its supernova remnant has been observed optically but was first detected at radio wavelengths. It is often known as 3C 10, a radio-source designation, although increasingly as Tycho's supernova remnant. Historic description The appearance of the Milky Way supernova of 1572 belongs among the most important observation events in the history of astronomy. The appearance of the "new star" helped to revise ancient models of the heavens and to speed on a revolution in astronomy that began with the realisation of the need to produce better astrometric star catalogues, and thus the need for more precise astronomical observing instruments. It also challenged the Aristotelian dogma of the unchangeability of the realm of stars. The supernova of 1572 is often called "Tycho's supernova", because of Tycho Brahe's extensive work De nova et nullius aevi memoria prius visa stella ("Concerning the Star, new and never before seen in the life or memory of anyone", published in 1573 with reprints overseen by Johannes Kepler in 1602 and 1610), a work containing both Brahe's own observations and the analysis of sightings from many other observers. Comparisons between Brahe's observations and those of Spanish scientist Jerónimo Muñoz revealed that the object was more distant than the Moon. This led Brahe to approach the Great Comet of 1577 as an astronomical body as well. Other Europeans to sight the supernova included Wolfgang Schuler, Christopher Clavius, Thomas Digges, John Dee, Francesco Maurolico, Tadeáš Hájek and . In England, Queen Elizabeth had the mathematician and astrologer Thomas Allen come and visit "to have his advice about the new star that appeared in the Swan or Cassiopeia ... to which he gave his judgement very learnedly", as the antiquary John Aubrey recorded in his memoranda a century later. In Ming dynasty China, the star became an issue between Zhang Juzheng and the young Wanli Emperor: in accordance with the cosmological tradition, the emperor was warned to consider his misbehavior, since the new star was interpreted as an evil omen. The more reliable contemporary reports state that the new star itself burst forth soon after November 2, 1572 and by November 11 it was already brighter than Jupiter. Around November 16, 1572, it reached its peak brightness at about magnitude −4.0, with some descriptions giving it as equal to Venus when that planet was at its brightest. Contrarily, Brahe described the supernova as "brighter than Venus". The supernova remained visible to the naked eye into early 1574, gradually fading until it disappeared from view. Supernova The supernova was classified as type I on the basis of its historical light curve soon after type I and type II supernovae were first defined on the basis of their spectra. The X-ray spectrum of the remnant showed that it was almost certainly of type Ia, but its detailed classification within the type Ia class continued to be debated until the spectrum of its light at peak luminosity was measured in a light echo in 2008. This gave final confirmation that it was a normal type Ia. The classification as a type Ia supernova of normal luminosity allows an accurate measure of the distance to SN 1572. The peak absolute magnitude can be calculated from the B-band decline rate to be . Given estimates of the peak apparent magnitude and the known extinction of magnitudes, the distance is kpc. Supernova remnant The distance to the supernova remnant has been estimated to between 2 and 5 kpc (approx. 6,500 and 16,300 light-years), with recent studies suggesting a narrower range of 2.5 and 3 kpc (approximately 8,000 and 9,800 light-years). Tycho's SNR has a roughly spherical morphology and spreads over an angular diameter of about 8 arcminutes. Its physical size corresponds to radius of the order of a few parsecs. Its measured expansion rate is about 11–12%/year in radio and X-ray. The average forward shock speed is between 4,000 and 5,000 km/s, dropping to lower speed when encountering local interstellar clouds. An older source says that the gas shell has reached an apparent diameter of 3.7 arcminutes. Initial radio detection The search for a supernova remnant was futile until 1952, when Robert Hanbury Brown and Cyril Hazard reported a radio detection at 158.5 MHz, obtained at the Jodrell Bank Observatory. This was confirmed, and its position more accurately measured in 1957 by Baldwin and Edge using the Cambridge Radio Telescope working at a wavelength of . The remnant was also identified tentatively in the second Cambridge Catalogue of Radio Sources as object "2C 34", and more firmly as "3C 10" in the third Cambridge list. There is no dispute that 3C 10 is the remnant of the supernova observed in 1572–1573. Following a 1964 review article by Minkowski, the designation 3C 10 appears to be that most commonly used in the literature when referring to the radio remnant of B Cas, although some authors use the tabulated galactic designation G120.7+2.1 and many authors commonly refer to it as Tycho's supernova remnant. Because the radio remnant was reported before the optical supernova-remnant wisps were discovered, the designation 3C 10 is used by some to signify the remnant at all wavelengths. X-ray observation An X-ray source designated Cepheus X-1 (or Cep X-1) was detected by the Uhuru X-ray observatory at 4U 0022+63. Earlier catalog designations are X120+2 and XRS 00224+638. Cepheus X-1 is actually in the constellation Cassiopeia, and it is SN 1572, the Tycho SNR. Optical detection The supernova remnant of B Cas was discovered in the 1960s by scientists with a Palomar Mountain telescope as a very faint nebula. It was later photographed by a telescope on the international ROSAT spacecraft. The supernova has been confirmed as Type Ia, in which a white dwarf star has accreted matter from a companion until it approaches the Chandrasekhar limit and explodes. This type of supernova does not typically create the spectacular nebula more typical of Type II supernovas, such as SN 1054 which created the Crab Nebula. A shell of gas is still expanding from its center at about 9,000 km/s. A recent study indicates a rate of expansion below 5,000 km/s. Companion star In October 2004, a letter in Nature reported the discovery of a G2 star, similar in type to our own Sun and named Tycho G. It is thought to be the companion star that contributed mass to the white dwarf that ultimately resulted in the supernova. A subsequent study, published in March 2005, revealed further details about this star: Tycho G was probably a main-sequence star or subgiant before the explosion, but some of its mass was stripped away and its outer layers were shock-heated by the supernova. Tycho G's current velocity is perhaps the strongest evidence that it was the companion star to the white dwarf, as it is traveling at a rate of 136 km/s, which is more than four times faster than the mean velocity of other stars in its stellar neighbourhood. This find has been challenged in recent years. The star is relatively far away from the center and does not show rotation which might be expected of a companion star. In Gaia DR2, the star was calculated to be light-years away, on the lower end of SN 1572's possible range of distances, which in turn lowered the calculated velocity from 136 km/s to only 56 km/s. In literature In the ninth episode of James Joyce's Ulysses, Stephen Dedalus associates the appearance of the supernova with the youthful William Shakespeare, and in the November 1998 issue of Sky & Telescope, three researchers from Southwest Texas State University, Don Olson and Russell Doescher of the Physics Department and Marilynn Olson of the English Department, argued that this supernova is described in Shakespeare's Hamlet, specifically by Bernardo in Act I, Scene i. The supernova inspired the poem "Al Aaraaf" by Edgar Allan Poe. The protagonist in Arthur C. Clarke's 1955 short story "The Star" casually mentions the supernova. It is a major element in Frederik Pohl's spoof science article, "The Martian Star-Gazers", first published in Galaxy Science Fiction Magazine in 1962.	Physical sciences	Notable transient events	Astronomy
2149536	https://en.wikipedia.org/wiki/Reptiliomorpha	Reptiliomorpha	Reptiliomorpha (meaning reptile-shaped; in PhyloCode known as Pan-Amniota) is a clade containing the amniotes and those tetrapods that share a more recent common ancestor with amniotes than with living amphibians (lissamphibians). It was defined by Michel Laurin (2001) and Vallin and Laurin (2004) as the largest clade that includes Homo sapiens, but not Ascaphus truei (tailed frog). Laurin and Reisz (2020) defined Pan-Amniota as the largest total clade containing Homo sapiens, but not Pipa pipa, Caecilia tentaculata, and Siren lacertina. The informal variant of the name, "reptiliomorphs", is also occasionally used to refer to stem-amniotes, i.e. a grade of reptile-like tetrapods that are more closely related to amniotes than they are to lissamphibians, but are not amniotes themselves; the name is used in this meaning e.g. by Ruta, Coates and Quicke (2003). An alternative name, "Anthracosauria", is also commonly used for the group, but is confusingly also used for a more primitive grade of reptiliomorphs (Embolomeri) by Benton. While both anthracosaurs and/or embolomeres are suggested to be reptiliomorphs closer to amniotes, some recent studies either retain them as amphibians or argue that their relationships are still ambiguous and are more likely to be stem-tetrapods. As the exact phylogenetic position of Lissamphibia within Tetrapoda remains uncertain, it also remains controversial which fossil tetrapods are more closely related to amniotes than to lissamphibians, and thus, which ones of them were reptiliomorphs in any meaning of the word. The two major hypotheses for lissamphibian origins are that they are either descendants of dissorophoid temnospondyls or microsaurian "lepospondyls". If the former (the "temnospondyl hypothesis") is true, then Reptiliomorpha includes all tetrapod groups that are closer to amniotes than to temnospondyls. These include the diadectomorphs, seymouriamorphs, most or all "lepospondyls", gephyrostegids, and possibly the embolomeres and chroniosuchians. In addition, several "anthracosaur" genera of uncertain taxonomic placement would also probably qualify as reptiliomorphs, including Solenodonsaurus, Eldeceeon, Silvanerpeton, and Casineria. However, if lissamphibians originated among the lepospondyls according to the "lepospondyl hypothesis", then Reptiliomorpha refers to groups that are closer to amniotes than to lepospondyls. Few non-amniote groups would count as reptiliomorphs under this definition, although the diadectomorphs are among those that qualify. Changing definitions The name Reptiliomorpha was coined by Professor Gunnar Säve-Söderbergh in 1934 to designate amniotes and various types of late Paleozoic tetrapods that were more closely related to amniotes than to living amphibians. In his view, the amphibians had evolved from fish twice, with one group composed of the ancestors of modern salamanders and the other, which Säve-Söderbergh referred to as Eutetrapoda, consisting of anurans (frogs), amniotes, and their ancestors, with the origin of caecilians being uncertain. Säve-Söderbergh's Eutetrapoda consisted of two sister-groups: Batrachomorpha, containing anurans and their ancestors, and Reptiliomorpha, containing anthracosaurs and amniotes. Säve-Söderbergh subsequently added Seymouriamorpha to his Reptiliomorpha as well. Alfred Sherwood Romer rejected Säve-Söderbergh's theory of a biphyletic amphibia and used the name Anthracosauria to describe the "labyrinthodont" lineage from which amniotes evolved. In 1970, the German paleontologist Alec Panchen took up Säve-Söderbergh's name for this group as having priority, but Romer's terminology is still in use, e.g. by Carroll (1988 and 2002) and by Hildebrand & Goslow (2001). Some writers preferring phylogenetic nomenclature use Anthracosauria. In 1956, Friedrich von Huene included both amphibians and anapsid reptiles in the Reptiliomorpha. This included the following orders: Anthracosauria, Seymouriamorpha, Microsauria, Diadectomorpha, Procolophonia, Pareiasauria, Captorhinidia, Testudinata. Michael Benton (2000, 2004) made it the sister-clade to Lepospondyli, containing "anthracosaurs" (in the strict sense, i.e. Embolomeri), seymouriamorphs, diadectomorphs and amniotes. Subsequently, Benton included lepospondyls in Reptiliomorpha as well. However, when considered in a Linnean framework, Reptiliomorpha is given the rank of superorder and includes only reptile-like amphibians, not their amniote descendants. Several phylogenetic studies indicate that amniotes and diadectomorphs share a more recent common ancestor with lepospondyls than with seymouriamorphs, Gephyrostegus and Embolomeri (e.g. Laurin and Reisz, 1997, 1999; Ruta, Coates and Quicke, 2003; Vallin and Laurin, 2004; Ruta and Coates, 2007). Lepospondyls are one of the groups of tetrapods suggested to be ancestors of living amphibians; as such, their potential close relationship to amniotes has important implications for the content of Reptiliomorpha. Assuming that lissamphibians aren't descended from lepospondyls but from a different group of tetrapods, e.g. from temnospondyls, it would mean that Lepospondyli belonged to Reptiliomorpha sensu Laurin (2001), as it would make them more closely related to amniotes than to lissamphibians. On the other hand, if lissamphibians are descended from lepospondyls, then not only Lepospondyli would have to be excluded from Reptiliomorpha, but seymouriamorphs, Gephyrostegus and Embolomeri would also have to be excluded from this group, as this would make them more distantly related to amniotes than living amphibians are. In that case, the clade Reptiliomorpha sensu Laurin would contain, apart from Amniota, only diadectomorphs and possibly also Solenodonsaurus. Characteristics Gephyrostegids, seymouriamorphs and diadectomorphs were land-based, reptile-like amphibians, while embolomeres were aquatic amphibians with long body and short limbs. Their anatomy falls between the mainly aquatic Devonian labyrinthodonts and the first reptiles. University of Bristol paleontologist Professor Michael J. Benton gives the following characteristics for the Reptiliomorpha (in which he includes embolomeres, seymouriamorphs and diadectomorphs): narrow premaxillae (less than half the skull width) vomers taper forward phalangeal formulae (number of joints in each toe) of foot 2.3.4.5.4–5 Cranial morphology The groups traditionally assigned to Reptiliomorpha, i.e. embolomeres, seymouriamorphs and diadectomorphs, differed from their contemporaries, the non-reptiliomorph temnospondyls, in having a deeper and taller skull, but retained the primitive kinesis (loose attachment) between the skull roof and the cheek (with exception of some specialized taxa, such as Seymouria, in which the cheek was solidly attached to the skull roof). The deeper skull allowed for laterally placed eyes, contrary to the dorsally placed eyes commonly found in amphibians. The skulls of the group are usually found with fine radiating grooves. The quadrate bone in the back of the skull held a deep otic notch, likely holding a spiracle rather than a tympanum. Postcranial skeleton The vertebrae showed the typical multi-element construction seen in labyrinthodonts. According to Benton, in the vertebrae of "anthracosaurs" (i.e. Embolomeri) the intercentrum and pleurocentrum may be of equal size, while in the vertebrae of seymouriamorphs the pleurocentrum is the dominant element and the intercentrum is reduced to a small wedge. The intercentrum gets further reduced in the vertebrae of amniotes, where it becomes a thin plate or disappears altogether. Unlike most labyrinthodonts, the body was moderately deep rather than flat, and the limbs were well-developed and ossified, indicating a predominantly terrestrial lifestyle except in secondarily aquatic groups. Each foot held five digits, the pattern seen in their amniote descendants. They did, however, lack the reptilian type of ankle bone that would have allowed the use of the feet as levers for propulsion rather than as holdfasts. Physiology The general build was heavy in all forms, though otherwise very similar to that of early reptiles. The skin, at least in the more advanced forms probably had a water-tight epidermal horny overlay, similar to the one seen in today's reptiles, though they lacked horny claws. In chroniosuchians and some seymouriamorphs, like Discosauriscus, dermal scales are found in post-metamorphic specimens, indicating they may have had a "knobbly", if not scaly, appearance. With reptiliomorph anthracosaurs having evolved small near-circular keratinous scales, their amniote descendants further covered almost their entire body with them, and also formed claws of keratin, with both scales and claws making cutaneous respiration and water absorption impossible, making them breathe through their mouths and nostrils, and drink water through mouth. Seymouriamorphs reproduced in amphibian fashion with aquatic eggs that hatched into larvae (tadpoles) with external gills; it is unknown how other tetrapods traditionally assigned to Reptiliomorpha reproduced. Evolutionary history Early reptiliomorphs During the Carboniferous and Permian periods, some tetrapods started to evolve towards a reptilian condition. Some of these tetrapods (e.g. Archeria, Eogyrinus) were elongate, eel-like aquatic forms with diminutive limbs, while others (e.g. Seymouria, Solenodonsaurus, Diadectes, Limnoscelis) were so reptile-like that until quite recently they actually had been considered to be true reptiles, and it is likely that to a modern observer they would have appeared as large to medium-sized, heavy-set lizards. Several groups however remained aquatic or semiaquatic. Some of the chroniosuchians show the build and presumably habits of modern crocodiles and were probably also similar to crocodylians in that they were river-side predators. While some other Chroniosuchians possessed elongated newt- or eel-like bodies. The two most terrestrially adapted groups were the medium-sized insectivorous or carnivorous Seymouriamorpha and the mainly herbivorous Diadectomorpha, with many large forms. The latter group has, in most analysis, the closest relatives of the Amniotes. From aquatic to terrestrial eggs Their terrestrial life style combined with the need to return to the water to lay eggs hatching to larvae (tadpoles) led to a drive to abandon the larval stage and aquatic eggs. A possible reason may have been competition for breeding ponds, to exploit drier environments with less access to open water, or to avoid predation on tadpoles by fish, a problem still plaguing modern amphibians. Whatever the reason, the drive led to internal fertilization and direct development (completing the tadpole stage within the egg). A striking parallel can be seen in the frog family Leptodactylidae, which has a very diverse reproductive system, including foam nests, non-feeding terrestrial tadpoles and direct development. The Diadectomorphans generally being large animals would have had correspondingly large eggs, unable to survive on land. Fully terrestrial life was achieved with the development of the amniote egg, where a number of membranous sacks protect the embryo and facilitate gas exchange between the egg and the atmosphere. The first to evolve was probably the allantois, a sack that develops from the gut/yolk-sack. This sack contains the embryo's nitrogenous waste (urea) during development, stopping it from poisoning the embryo. A very small allantois is found in modern amphibians. Later came the amnion surrounding the fetus proper, and the chorion, encompassing the amnion, allantois, and yolk-sack. Origin of amniotes Exactly where the border between reptile-like amphibians (non-amniote reptiliomorphs) and amniotes lies will probably never be known, as the reproductive structures involved fossilize poorly, but various small, advanced reptiliomorphs have been suggested as the first true amniotes, including Solenodonsaurus, Casineria and Westlothiana. Such small animals laid small eggs, 1 cm in diameter or less. Small eggs would have a small enough volume to surface ratio to be able to develop on land without the amnion and chorion actively affecting gas exchange, setting the stage for the evolution of true amniotic eggs. Although the first true amniotes probably appeared as early as the Middle Mississippian sub-epoch, non-amniote (or amphibian) reptiliomorph lineages coexisted alongside their amniote descendants for many millions of years. By the middle Permian the non-amniote terrestrial forms had died out, but several aquatic non-amniote groups continued to the end of the Permian, and in the case of the chroniosuchians survived the end Permian mass extinction, only to die out prior to the end of the Triassic. Meanwhile, the single most successful daughter-clade of the reptiliomorphs, the amniotes, continued to flourish and evolve into a staggering diversity of tetrapods including mammals, reptiles, and birds. Gallery	Biology and health sciences	Reptiliomorphs	Animals
2150441	https://en.wikipedia.org/wiki/Quantifier%20elimination	Quantifier elimination	Quantifier elimination is a concept of simplification used in mathematical logic, model theory, and theoretical computer science. Informally, a quantified statement " such that " can be viewed as a question "When is there an such that ?", and the statement without quantifiers can be viewed as the answer to that question. One way of classifying formulas is by the amount of quantification. Formulas with less depth of quantifier alternation are thought of as being simpler, with the quantifier-free formulas as the simplest. A theory has quantifier elimination if for every formula , there exists another formula without quantifiers that is equivalent to it (modulo this theory). Examples An example from mathematics says that a single-variable quadratic polynomial has a real root if and only if its discriminant is non-negative: Here the sentence on the left-hand side involves a quantifier , whereas the equivalent sentence on the right does not. Examples of theories that have been shown decidable using quantifier elimination are Presburger arithmetic, algebraically closed fields, real closed fields, atomless Boolean algebras, term algebras, dense linear orders, abelian groups, random graphs, as well as many of their combinations such as Boolean algebra with Presburger arithmetic, and term algebras with queues. Quantifier eliminator for the theory of the real numbers as an ordered additive group is Fourier–Motzkin elimination; for the theory of the field of real numbers it is the Tarski–Seidenberg theorem. Quantifier elimination can also be used to show that "combining" decidable theories leads to new decidable theories (see Feferman–Vaught theorem). Algorithms and decidability If a theory has quantifier elimination, then a specific question can be addressed: Is there a method of determining for each ? If there is such a method we call it a quantifier elimination algorithm. If there is such an algorithm, then decidability for the theory reduces to deciding the truth of the quantifier-free sentences. Quantifier-free sentences have no variables, so their validity in a given theory can often be computed, which enables the use of quantifier elimination algorithms to decide validity of sentences. Related concepts Various model-theoretic ideas are related to quantifier elimination, and there are various equivalent conditions. Every first-order theory with quantifier elimination is model complete. Conversely, a model-complete theory, whose theory of universal consequences has the amalgamation property, has quantifier elimination. The models of the theory of the universal consequences of a theory are precisely the substructures of the models of . The theory of linear orders does not have quantifier elimination. However the theory of its universal consequences has the amalgamation property. Basic ideas To show constructively that a theory has quantifier elimination, it suffices to show that we can eliminate an existential quantifier applied to a conjunction of literals, that is, show that each formula of the form: where each is a literal, is equivalent to a quantifier-free formula. Indeed, suppose we know how to eliminate quantifiers from conjunctions of literals, then if is a quantifier-free formula, we can write it in disjunctive normal form and use the fact that is equivalent to Finally, to eliminate a universal quantifier where is quantifier-free, we transform into disjunctive normal form, and use the fact that is equivalent to Relationship with decidability In early model theory, quantifier elimination was used to demonstrate that various theories possess properties like decidability and completeness. A common technique was to show first that a theory admits elimination of quantifiers and thereafter prove decidability or completeness by considering only the quantifier-free formulas. This technique can be used to show that Presburger arithmetic is decidable. Theories could be decidable yet not admit quantifier elimination. Strictly speaking, the theory of the additive natural numbers did not admit quantifier elimination, but it was an expansion of the additive natural numbers that was shown to be decidable. Whenever a theory is decidable, and the language of its valid formulas is countable, it is possible to extend the theory with countably many relations to have quantifier elimination (for example, one can introduce, for each formula of the theory, a relation symbol that relates the free variables of the formula). Example: Nullstellensatz for algebraically closed fields and for differentially closed fields.	Mathematics	Model theory	null
2150513	https://en.wikipedia.org/wiki/Green%20sea%20turtle	Green sea turtle	The green sea turtle (Chelonia mydas), also known as the green turtle, black (sea) turtle or Pacific green turtle, is a species of large sea turtle of the family Cheloniidae. It is the only species in the genus Chelonia. Its range extends throughout tropical and subtropical seas around the world, with two distinct populations in the Atlantic and Pacific Oceans, but it is also found in the Indian Ocean. The common name refers to the usually green fat found beneath its carapace, due to its diet strictly being seagrass, not to the color of its carapace, which is olive to black. The dorsoventrally flattened body of C. mydas is covered by a large, teardrop-shaped carapace; it has a pair of large, paddle-like flippers. It is usually lightly colored, although in the eastern Pacific populations, parts of the carapace can be almost black. Unlike other members of its family, such as the hawksbill sea turtle, C. mydas is mostly herbivorous. The adults usually inhabit shallow lagoons, feeding mostly on various species of seagrasses. The turtles bite off the tips of the blades of seagrass, which keeps the grass healthy. Like other sea turtles, green sea turtles migrate long distances between feeding grounds and hatching beaches. Many islands worldwide are known as Turtle Island due to green sea turtles nesting on their beaches. Females crawl out on beaches, dig nests, and lay eggs during the night. Later, hatchlings emerge, and scramble into the water. Those that reach maturity may live to 90 years in the wild. C. mydas is listed as endangered by the IUCN and CITES and is protected from exploitation in most countries. It is illegal to collect, harm, or kill them. In addition, many countries have laws and ordinances to protect nesting areas. However, turtles are still in danger due to human activity. In some countries, turtles and their eggs are still hunted for food. Pollution indirectly harms turtles at both population and individual scales. Many turtles die after being caught in fishing nets. In addition, real estate development often causes habitat loss by eliminating nesting beaches. Taxonomy The green sea turtle is a member of the tribe Cheloniini. A 1993 study clarified the status of genus Chelonia with respect to the other marine turtles. The carnivorous Eretmochelys (hawksbill), Caretta (loggerhead) and Lepidochelys (ridley) were assigned to the tribe Carettini. Herbivorous Chelonia warranted their status as a genus, while Natator (flatback) was further removed from the other genera than previously believed. The species was originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae as Testudo mydas. In 1868, Marie Firmin Bocourt named a particular species of sea turtle Chelonia agassizii, in honor of Swiss-American zoologist Louis Agassiz. This "species" was referred to as the "black sea turtle". Later research determined Bocourt's "black sea turtle" was not genetically distinct from C. mydas, and thus taxonomically not a separate species. These two "species" were then united as Chelonia mydas and populations were given subspecies status: C. mydas mydas referred to the originally described population, while C. mydas agassizi referred only to the Pacific population known as the Galápagos green turtle. This subdivision was later determined to be invalid and all species members were then designated Chelonia mydas. The oft-mentioned name C. agassizi remains an invalid junior synonym of C. mydas. The species' common name does not derive from any particular green external coloration of the turtle. Its name comes from the greenish color of the turtles' fat, which is only found in a layer between their inner organs and their shell. As a species found worldwide, the green turtle has many local names. In the Hawaiian language it is called honu, and it is locally known as a symbol of good luck and longevity. Description Its appearance is that of a typical sea turtle. C. mydas has a dorsoventrally flattened body, a beaked head at the end of a short neck, and paddle-like arms well-adapted for swimming. Adult green turtles grow to long. The average weight of mature individuals is and the average carapace length is . They are considered the second largest sea turtle in the United States, after the leatherback sea turtle. Exceptional specimens can weigh or even more, with the largest known C. mydas having weighed and measured in carapace length. Anatomically, a few characteristics distinguish the green turtle from the other members of its family. Unlike its close relative the hawksbill turtle, the green turtle's snout is very short and its beak is unhooked. The neck cannot be pulled into the shell. The sheath of the turtle's upper jaw possesses a denticulated edge, while its lower jaw has stronger, serrated, more defined denticulation. The dorsal surface of the turtle's head has a single pair of prefrontal scales. Its carapace is composed of five central scutes flanked by four pairs of lateral scutes. Underneath, the green turtle has four pairs of inframarginal scutes covering the area between the turtle's plastron and its shell. Mature C. mydas front appendages have only a single claw (as opposed to the hawksbill two), although a second claw is sometimes prominent in young specimens. The carapace of the turtle has various color patterns that change over time. Hatchlings of Chelonia mydas, like those of other marine turtles, have mostly black carapaces and light-colored plastrons. Carapaces of juveniles turn dark brown to olive, while those of mature adults are either entirely brown, spotted or marbled with variegated rays. Underneath, the turtle's plastron is hued yellow. C. mydas limbs are dark-colored and lined with yellow, and are usually marked with a large dark brown spot in the center of each appendage. Distribution The range of the green sea turtle extends throughout tropical and subtropical oceans worldwide. The two major subpopulations are the Atlantic and the eastern Pacific subpopulations. Each population is genetically distinct, with its own set of nesting and feeding grounds within the population's known range. One of the genetic differences between the two subpopulations is the type of mitochondrial DNA found in individual's cells. Individuals from rookeries in the Atlantic Ocean and Mediterranean Sea have a similar type of mitochondrial DNA, and individuals from the Pacific and Indian Oceans have another type of mitochondrial DNA. Their native range includes tropical to subtropical waters along continental coasts and islands between 30°N and 30°S. Since green sea turtles are a migrating species, their global distribution spans into the open ocean. The differences in mitochondrial DNA more than likely stems from the populations being isolated from each other by the southern tips of both South America and Africa with no warm waters for the green sea turtles to migrate through. The green sea turtle is estimated to inhabit coastal areas of more than 140 countries, with nesting sites in over 80 countries worldwide throughout the year. In the United States Atlantic coast, green sea turtles can be found from Texas and north to Massachusetts. In the United States Pacific coast, they have been found from southern California north to the southernmost tip of Alaska. The largest populations of green sea turtles within the United States coastline are in the Hawaiian Islands and Florida. Globally, the largest populations of sea turtles are in the Great Barrier Reef in Australia, and the Caribbean Sea. Atlantic subpopulation The green sea turtle can generally be found throughout the Atlantic Ocean. Although the species is most abundant in tropical climates, green sea turtles can also be found in temperate climates, and individuals have been spotted as far north as Canada in the western Atlantic, and the Cimbrian peninsular in the east. The subpopulation's southern range is known until past the southern tip of Africa in the east and Argentina in the western Atlantic. The major nesting sites can be found on various islands in the Caribbean, along the Atlantic coast of Florida in the United States, the eastern coast of the South American continent and most notably, on isolated North Atlantic islands. In the Caribbean, major nesting sites have been identified on Aves Island, the U.S. Virgin Islands, Puerto Rico, the Dominican Republic, and Costa Rica. In recent years, there are signs of increased nesting in the Cayman Islands. One of the region's most important nesting grounds is in Tortuguero in Costa Rica. In fact, the majority of the Caribbean region's C. mydas population hails from a few beaches in Tortuguero. Within United States waters, minor nesting sites have been noted in the states of Georgia, North Carolina, and South Carolina. Florida's east coast is the largest nesting site in the United States. Hutchinson Island in particular is a major nesting area in Florida waters. Florida has several annual nesting periods when local beaches are closed or cordoned off to protect nesting sites. According to Green Sea Turtle Watch, in 2015 more than 37,000 green sea turtle nests were documented in Florida, a record number. In addition to sporadic distribution of nesting sites, feeding grounds are much more widely distributed throughout Florida. Important feeding grounds in Florida include Indian River Lagoon, the Florida Keys, Florida Bay, Homosassa, Crystal River, and Cedar Key. Notable locations in South America include secluded beaches in Suriname and French Guiana. In the Southern Atlantic Ocean, the most notable nesting grounds for Chelonia mydas are found on the island of Ascension, hosts 6,000–13,000 turtle nests. Indo-Pacific subpopulation In the Pacific, its range reaches as far north as the southern coast of Alaska and as far south as Chile in the east. The turtle's distribution in the western Pacific reaches north to Japan and southern parts of Russia's Pacific coast, and as far south as the northern tip of New Zealand and a few islands south of Tasmania. Significant nesting grounds are scattered throughout the entire Pacific region, including Mexico, the Hawaiian Islands, the South Pacific, the northern coast of Australia, and Southeast Asia. Major Indian Ocean nesting colonies include India, Pakistan, Sri Lanka and other coastal countries. The turtles can also be found throughout the Indian Ocean; the east coast of the African continent hosts a few nesting grounds, including islands in the waters around Madagascar. Specific nesting grounds Nesting grounds are found all along the Mexican coast. These turtles feed in seagrass pastures in the Gulf of California. Green turtles belonging to the distinct Hawaiian subpopulation nest at the protected French Frigate Shoals some west of the Hawaiian Islands. In the Philippines, green turtles nest in the Turtle Islands along with closely related hawksbill turtles. In December 2007, fishermen using a hulbot-hulbot (a type of fish net) accidentally caught an , long and wide turtle off Barangay Bolong, Zamboanga City, Philippines. December is breeding season near the Bolong beach. An annual presence is recorded in the Gulf of Panama, on the Isla Parida island. Local activists also moving some turtle nests to the coast, in the vicinity of the small town of Malena, to save and increase the turtle population in the safe place. Indonesia has a few nesting beaches, one in the Meru Betiri National Reserve in East Java. Off the north-eastern and northern coasts of Australia, the Great Barrier Reef has two genetically distinct populations; one north and one south. Within the reef, 20 separate locations consisting of small islands and cays were identified as nesting sites for either population of C. mydas. Of these, the most important is on Raine Island. In the Torres Strait there is a large rookery on Bramble Cay. The Coral Sea has nesting areas of world significance. Major nesting sites are common on either side of the Arabian Sea, both in Ash Sharqiyah, Oman, and along the coast of Karachi, Pakistan. Some specific beaches there, such as Hawke's Bay and Sandspit, are common to both C. mydas and L. olivacea subpopulation. Sandy beaches along Sindh and Balochistan are nesting sites. Some off the Pakistani coast, Astola island is another nesting beach. Galápagos green turtle The population that has often been known as the Galápagos green turtle have been recorded and observed in the Galápagos as far back as the 17th century by William Dampier. Not much attention has been paid to them due to the overwhelming research done on the Galápagos giant tortoises. Only over the last 30 years have extensive studies been performed covering the behaviors of the Galápagos green turtles. Much of the debate that has surrounded them recently is over the binomial classification of the species. At one point the name Chelonia agassizii was applied to this population as a separate species. Analysis of mitochondrial and nuclear DNA of 15 nesting beaches, however, has demonstrated that there is not only no significant distinction of this population but that it would be paraphyletic to recognise it. As such the species name Chelonia agassizzii is considered a junior synonym of Chelonia mydas as such it is considered as a local variant of the populations of the East Pacific waters and those of other nesting areas. The morphological distinctiveness of the Galápagos green turtle has given rise to the debate, but evidence of taxonomic distinctiveness is best served using the combination of multiple datasets. The two most notable morphological distinctions are the considerably smaller adult size and the much darker pigmentation of the carapace, plastron, and extremities. Other distinctions are the curving of the carapace above each hind flipper, the more dome-shaped carapace, and the very long tail of adult males. Three possibilities have arisen from their unique characteristics: agassizii is a separate species from C. mydas, it is a subspecies of green sea turtle, or it is simply a color mutation. These facts have led to the debate over binomial separation however due to the significance of the DNA testing results there have been no distinctions made at this time. At a meeting for sea turtle scientists and their collaborators in 2000, the evidence for the taxonomic position of the Galápagos green turtle was reviewed and a majority among the participants supported treating it as a population or subspecies of the green turtle (instead of a separate species). However, this is possibly a case of political taxonomy. As such the three major international checklists that cover turtles of the world Reptile Database the checklist of Fritz and Havas (2007) and the IUCN Checklist (TTWG 2017) all consider this a junior synonym. Habitat Green sea turtles move across three habitat types, depending on their life stage. They lay eggs on beaches. Mature turtles spend most of their time in shallow, coastal waters with lush seagrass beds. Adults frequent inshore bays, lagoons, and shoals with lush seagrass meadows. Entire generations often migrate between one pair of feeding and nesting areas. Green sea turtles, Chelonia mydas, are classified as an aquatic species and are distributed around the globe in warm tropical to subtropical waters. The environmental parameter that limits the distribution of the turtles is ocean temperatures below 7 to 10 degrees Celsius. Within their geographical range, the green sea turtles generally stay near continental and island coastlines. Near the coastlines, the green sea turtles live within shallow bays and protected shores. In these protected shores and bays, the green sea turtle habitats include coral reefs, salt marshes, and nearshore seagrass beds. The coral reefs provide red, brown, and green algae for their diet and give protection from predators and rough storms within the ocean. The salt marshes and seagrass beds contain seaweed and grass vegetation, allowing ample habitat for the sea turtles. Turtles spend most of their first five years in convergence zones within the bare open ocean that surround them. These young turtles are rarely seen as they swim in deep, pelagic waters. Green sea turtles typically swim at . Ecology and behavior As one of the first sea turtle species studied, much of what is known of sea turtle ecology comes from studies of green turtles. The ecology of C. mydas changes drastically with each stage of its life history. Newly emerged hatchlings are carnivorous, pelagic organisms, part of the open ocean mininekton. In contrast, immature juveniles and adults are commonly found in seagrass meadows closer inshore as herbivorous grazers. Diet The diet of green turtles changes with age. Juveniles are carnivorous, but as they mature they become omnivorous. Young sea turtles eat fish and their eggs, sea hare eggs, hydrozoans, bryozoans, molluscs, jellyfish, small invertebrates, echinoderms, tunicates, insects, worms, sponges, algae, sea grasses, leaves, tree bark, and crustaceans. Green sea turtles have a relatively slow growth rate because of the low nutritional value of their diet. Body fat turns green because of the consumed vegetation. This diet shift has an effect on the green turtle's skull morphology. Their serrated jaw helps them chew green and red algae (such as filamentous red alga (Bostrychia), red moss (Caloglossa), freshwater red algae (Compsopogon), lobster horns (Polysiphonia), sea lettuce (Ulva lactuca), green seaweed (Gayralia), and crinkle grass (Rhizoclonium)) and sea grasses. They also consume large quantities of wetland plants such as Avicennia schaueriana and Sporobolus alterniflorus, which are commonly found in salt marshes. Most adult sea turtles are strictly herbivorous. Predators and parasites Only some human beings and the larger sharks feed on C. mydas adults. Specifically, tiger sharks (Galeocerdo cuvier) hunt adults in Hawaiian waters. The tiger shark is the main predator of the green turtle as it will prey on green turtles of all sizes. The tiger shark has often been seen feeding on green turtles near their nesting beaches because they are restricted in the area of their nesting beaches and vulnerable to predation. Juveniles and new hatchlings have significantly more predators, including crabs, small marine mammals and shorebirds. Additionally, their eggs are vulnerable to predation by scavengers like red foxes and golden jackals. Green sea turtles have a variety of parasites including barnacles, leeches, protozoans, cestodes, and nematodes. Barnacles attach to the carapace, and leeches to the flippers and skin of the turtles, causing damage to the soft tissues and leading to blood loss. Protozoans, cestodes and nematodes lead to many turtle deaths because of the infections in the liver and intestinal tract they cause. The greatest disease threat to the turtle population is fibropapilloma, which produces lethal tumor growth on scales, lungs, stomach, and kidneys. Fibropapilloma is caused by a herpesvirus that is transmitted by leeches such as Ozobranchus branchiatus, a species of leech which feeds almost entirely on green sea turtles. Life cycle Green sea turtles migrate long distances between feeding sites and nesting sites; some swim more than to reach their spawning grounds. Beaches in Southeast Asia, India, islands in the western Pacific, and Central America are where green sea turtles breed. Mature turtles often return to the exact beach from which they hatched. Females usually mate every two to four years. Males, on the other hand, visit the breeding areas every year, attempting to mate. Mating seasons vary between populations. For most C. mydas in the Caribbean, mating season is from June to September. The French Guiana nesting subpopulation nests from March to June. In the tropics, green turtles nest throughout the year, although some subpopulations prefer particular times of the year. In Pakistan, Indian Ocean turtles nest year-round, but prefer the months of July to December. Sea turtles return to the beaches on which they were born to lay their own eggs. The reason for returning to native beaches may be that it guarantees the turtles an environment that has the necessary components for their nesting to be successful. These include a sandy beach, easy access for the hatchlings to get to the ocean, the right incubation temperatures, and low probability of predators that may feed on their eggs. Over time these turtles have evolved these tendencies to return to an area that has provided reproductive success for many generations. Their ability to return to their birthplace is known as natal homing. The males also return to their birthplaces in order to mate. These males that return to their homes know they will be able to find mates because the females born there also return to breed. By doing this, the green sea turtles are able to improve their reproductive success and is why they are willing to expend the energy to travel thousands of miles across the ocean in order to reproduce. Mating behaviour is similar to other marine turtles. Female turtles control the process. A few populations practice polyandry, although this does not seem to benefit hatchlings. After mating in the water, the female moves above the beach's high tide line, where she digs a hole in depth with her hind flippers and deposits her eggs. The hole is then covered up again. Clutch size ranges between 85 and 200, depending on the age of the female. This process takes about an hour to an hour and a half. After the nest is completely covered, she returns to the sea. The female will do this 3 to 5 times in one season. The eggs are round and white, and about in diameter. The hatchlings remain buried for days until they all emerge together at night. The temperature of the nest determines the sex of the turtles at around the 20–40 day mark. Green Sea Turtles are type 1a, meaning males develop at cooler temperatures while females are produced under hot temperatures. At around 50 to 70 days, the eggs hatch during the night, and the hatchlings instinctively head directly into the water. This is the most dangerous time in a turtle's life. As they walk, predators, such as gulls and crabs, feed on them. A significant percentage never make it to the ocean. Little is known of the initial life history of newly hatched sea turtles. Juveniles spend three to five years in the open ocean before they settle as still-immature juveniles into their permanent shallow-water lifestyle. It is speculated that they take twenty to fifty years to reach sexual maturity. Individuals live up to eighty years in the wild. It is estimated that only 1% of hatchlings reach sexual maturity. Each year on Ascension Island in the South Atlantic, C. mydas females create 6,000 to 25,000 nests. They are among the largest green turtles in the world; many are more than in length and weigh up to . Breathing and sleep Sea turtles spend almost all their lives submerged, but must breathe air for the oxygen needed to meet the demands of vigorous activity. With a single explosive exhalation and rapid inhalation, sea turtles can quickly replace the air in their lungs. The lungs permit a rapid exchange of oxygen and prevent gases from being trapped during deep dives. Sea turtle blood can deliver oxygen efficiently to body tissues even at the pressures encountered during diving. During routine activity, green and loggerhead turtles dive for about four to five minutes, and surface to breathe for one to three seconds. Turtles can rest or sleep underwater for several hours at a time, but submergence time is much shorter while diving for food or to escape predators. Breath-holding ability is affected by activity and stress, which is why turtles quickly drown in shrimp trawlers and other fishing gear. During the night while sleeping and to protect themselves from potential predators, the adults wedge themselves under rocks below the surface and under ledges in reefs and coastal rocks. Many green sea turtles have been observed in returning to the same sleeping location night after night. Physiology and sensory modalities Green sea turtles tend to have good vision, well adapted to a life at sea. The turtles can see many colors, but are most sensitive to light from violet to yellow or wavelengths of 400 to 600 nanometers. They do not see many colors in the orange to red portion of the light spectrum. On land, however, the sea turtles are nearsighted because the lenses in the eyes are spherical and adjusted to refraction underwater. Sea turtles have no external ear and only one ear bone, called the columella. With one ear bone, the turtles can hear only low frequency sounds, from 200 to 700 Hz. Sounds can also be detected through vibrations of the head, backbone, and shell. The nose of the turtle has two external openings and connects to the roof of the mouth through internal openings. The lower surface of the nasal passage has two sets of sensory cells called the Jacobson's organ. The turtle can use this organ to smell by pumping water in and out of its nose. Since green sea turtles migrate long distances during breeding seasons, they have special adaptive systems in order to navigate. In the open ocean, the turtles navigate using wave directions, sun light, and temperatures. The sea turtles also contain an internal magnetic compass. They can detect magnetic information by using magnetic forces acting on the magnetic crystals in their brains. Through these crystals, they can sense the intensity of Earth's magnetic field and are able to make their way back to their nesting grounds or preferred feeding grounds. Natal homing is an animal's ability to return to its birthplace in order to reproduce. Natal homing is found in all species of sea turtles and in other animals such as salmon. How these turtles are able to return to their birthplace is an interesting phenomenon. Many researchers believe that sea turtles use a process called imprinting, which is a special type of learning that occurs when turtles first hatch that allows them to recognize their native beach. There are two types of imprinting that are thought to be the reason turtles can find these beaches. The first is the chemical imprinting hypothesis. This hypothesis states that much like salmon, sea turtles are able to use olfactory cues and senses to smell their way home. However, a problem with this hypothesis is that some turtles travel thousands of miles to return to their native beaches, and the scents from that area are not likely to travel and be distinguishable from that distance. The second hypothesis is the geomagnetic. This hypothesis states that as it hatches, a young turtle will imprint on the magnetic field of the beach they are born on. This hypothesis strongly correlates to the method which sea turtles use to navigate the earth. In order to tolerate the constant heat loss in the water, sea turtles have the ability to shunt blood away from tissues that are tolerant of low oxygen levels toward the heart, brain, and central nervous system. Other mechanisms include basking on warm beaches and producing heat through their activity and movements of their muscles. Basking turtles sometimes look like they are crying because behind the turtles eye is the lachrymal gland which stores excess salt from the sea water, which then expels through the turtles eye. In the winter months, turtles living at higher latitudes can hibernate for a short period in the mud. Unique characteristics and features Green sea turtles can reach up to 40 miles per hour when swimming, making them the fastest sea turtle. The green sea turtles exhibit sex differences by their development and appearance. As adult turtles, males are easily distinguishable from the females by having a longer tail (visibly extending past the shell) and longer claws on the front flippers. The hatching time and sex of the turtles are determined by the incubation temperature of the nest. Hatchings occur more quickly in nests that are warmer than nests that are in cooler conditions. Warm nesting sites above 30 degrees Celsius favor the development of females, whereas nesting sites below 30 degrees Celsius produce males. The position of the egg in the nest also affects sex-determination. Eggs in the center tend to hatch as females due to the warmer conditions within the nest. Green sea turtles play an essential role in the ecosystem in which they live. In the seagrass beds, the turtles feed on the seagrass by trimming only the top and leaving the roots of the plant. Through their feeding technique, the turtles help to improve the health and growth of the seagrass beds. The healthy seagrass beds that the turtles provide give habitat and feeding grounds for many species of fish and crustaceans. On the nesting beaches, the green sea turtles provide key nutrients for the ecosystem through their hatched egg shells. In their coral reef habitat, the green sea turtles have a symbiotic interaction with reef fish, including the yellow tang. The yellow tang fish swims along with the turtle and feeds on the algae, barnacles, and parasites on its shell and flippers. This species interaction provides food for the yellow tang and provides a necessary cleaning and smoothing of the turtle's shell. This cleaning helps the turtles swim by reducing the amount of drag and improves their health. Importance to humans Historically, the turtles' skin was tanned and used to make handbags, especially in Hawaii. Ancient Chinese considered the flesh of sea turtles a culinary delicacy, including and especially C. mydas. Particularly for this species, the turtle's fat, cartilage, and flesh, known as calipee, are sought as ingredients for making turtle soup, a popular 19th-century English and American dish. In Java, Indonesia, sea turtle eggs were a popular delicacy. However, the turtle's flesh is regarded as ḥarām or "unclean" under Islamic law (Islam is Java's primary religion). In Bali, turtle meat was a prominent feature at ceremonial and religious feasts. Turtles were harvested in the remotest parts of the Indonesian archipelago. Bali has been importing sea turtles since the 1950s, as its own turtle supplies became depleted. The mostly Hindu Balinese do not eat the eggs, but sell them instead to local Muslims. Commercial farms, such as the Cayman Turtle Farm in the West Indies, once bred them for commercial sale of turtle meat, turtle oil (rendered from the fat), turtle shell, and turtle leather made from the skin. The farm's initial stock was in large part from "doomed" eggs removed from nests threatened by erosion, flooding, or in chemically hostile soil. The farms held as many as 100,000 turtles at any one time. When the international markets were closed by regulations that did not allow even farm-bred turtle products to be exported internationally, the surviving farm became primarily a tourist attraction, supporting 11,000 turtles. Initially started as Mariculture Ltd., then Cayman Turtle Farm Ltd and subsequently branded Boatswain's Beach, in 2010 the farm's brandname was changed to Cayman Turtle Farm: Island Wildlife Encounter. Sea turtles are integral to the history and culture of the Cayman Islands. When the islands were first discovered by Christopher Columbus in 1503, he named them "Las Tortugas" because of the abundance of sea turtles in the waters around the islands. Many of the earliest visitors came to the Cayman Islands to capture the turtles as a source of fresh meat during long voyages. The green turtle is a national symbol displayed as part of the coat of arms of the Cayman Islands, which also forms part of the national flag of the Cayman Islands. The country's currency uses a turtle as the watermark in its banknotes. A stylised sea turtle nicknamed "Sir Turtle" is the mascot of the national airline Cayman Airways and is part of the livery of its aircraft. A ki'i pōhaku (petroglyph) of a green sea turtle (or honu) can be found on the Big Island of Hawaii in the Pu'u Loa lava fields. The green sea turtle has always held a special meaning for Hawaiians and this petroglyph shows its importance; it may date to when the Hawaiian Islands first became populated. The turtle symbolizes a navigator that can find his way home time after time. This symbol mirrors the real life of the green Hawaiian turtle as it will swim hundreds of miles to lay its eggs at its own place of birth. Though there are other myths as well, some Hawaiian legends say the honu were the first to guide the Polynesians to the Hawaiian Islands. Hawaiians revere the turtle and the legend of Kailua, a turtle who could take the form of a girl at will. In human form, she looked after the children playing on Punalu'u Beach. Conservation In recent decades, sea turtles have moved from unrestricted exploitation to global protection, with individual countries providing additional protection, although serious threats remain unabated. All populations are considered "threatened". Threats Human action presents both intentional and unintentional threats to the species' survival. Intentional threats include continued hunting, poaching and egg harvesting. More dangerous are unintentional threats, including boat strikes, fishermen's nets that lack turtle excluder devices, pollution and habitat destruction. Chemical pollution may create tumors; effluent from harbors near nesting sites may create disturbances; and light pollution may disorient hatchlings. With chemical pollution present, there is a development of tar balls that is often eaten by green sea turtles in a confusion of their food. Tar balls cause the green sea turtle to ingest toxins that can block their guts and cause swelling of the tissue, displacing the liver and intestines. Habitat loss usually occurs due to human development of nesting areas. Beach-front construction, land "reclamation" and increased tourism are examples of such development. An infectious tumor-causing disease, fibropapillomatosis, is also a problem in some populations. The disease kills a sizeable fraction of those it infects, though some individuals seem to resist the disease. In addition, at least in the Southwestern Atlantic (Río de la Plata, Uruguay), exotic invasive species such as the rapa whelk Rapana venosa, were reported massively bio-fouling immature green turtles, reducing buoyancy, increasing drag, and causing severe injuries to the carapace. Because of these threats, many populations are in a vulnerable state. Pacific green turtles' foraging habitats are poorly understood and mostly unknown. These foraging grounds are most likely along the coast of Baja California, Mexico and southern California, in which these turtles have a high risk of incidental capture by coastal fisheries. The main mortality factor for these turtles is the shrimp trawlers in Mexico, in which many of these turtles go undocumented. The only foraging area that has been identified is San Diego Bay, but it is heavily polluted with metals and PCBs. These contaminants have a negative effect on the ocean environment, and have been shown to cause lesions and sometimes mortality. Green turtles also are threatened by entanglement and ingestion of plastic. In San Diego Bay, an adult green turtle was found dead with monofilament netting tightly packed in its esophagus. In addition there are indications that global climate change is affecting the ability of green turtle populations in Australia to produce males due to their temperature-dependent sex determination and the rising temperatures in the northern Great Barrier Reef region. Construction of new thermal power stations can raise local water temperature, which is also said to be a threat. Green sea turtles are the most commonly traded species along Java's south coast and are sold in the form of whole, stuffed animals or turtle oil, locally known as "minyak bulus". The geographer James J. Parsons' book titled The Green Turtle and Man played a special role in the conservation movement to save the species from extinction. Global initiatives The International Union for Conservation of Nature (IUCN) has repeatedly listed green sea turtles in its Red List under differing criteria. In 1982, they officially classified it as an endangered species. The 1986, 1988, 1990, 1994, and the landmark 1996 edition of the IUCN Red List, retained the listing. In 2001, Nicholas Mrosovsky filed a delisting petition, claiming some green turtle populations were large, stable and in some cases, increasing. At the time, the species was listed under the strict EN A1abd criteria. The IUCN Standards and Petitions Subcommittee ruled that visual counts of nesting females could not be considered "direct observation" and thus downgraded the species' status to EN A1bd—retaining the turtle's endangered status. In 2004, the IUCN reclassified C. mydas as endangered under the EN A2bd criteria, which essentially states the wild populations face a high risk of extinction because of several factors. These factors include a probable population reduction of more than 50% over the past decade as estimated from abundance indices and by projecting exploitation levels. On 3 May 2007, C. mydas was listed on Appendix I of the Convention on International Trade in Endangered Species (CITES) as a member of the family Cheloniidae. The species was originally listed on Appendix II in 1975. The entire family was moved to Appendix I in 1977, with the exception of the Australian population of C. mydas. In 1981, the Australian population joined the rest. The Appendix I listing prohibits commercial international trade in the species (including parts and derivatives). The Zoological Society of London has listed the reptile as an EDGE species. The Mediterranean population is listed as critically endangered. The eastern Pacific, Hawaiian and Southern California subpopulations are designated threatened. Specific Mexican subpopulations are listed as endangered. The Florida population is listed as endangered. The World Wide Fund for Nature has labeled populations in Pakistan as "rare and declining". Since 1999, the Florida Aquarium has led extensive sea turtle rehabilitation efforts and visitor and community education & conservation platforms to advance sea turtle protection. Over a 20-year period, the aquarium received 200 sea turtles, and while not all could be released due to the nature of their injuries or illnesses, 180 were successfully released. In 2019, they opened a state-of-the-art Sea Turtle Rehabilitation Center in Apollo Beach, Florida. In the first year, The Florida Aquarium Animal Response Team managed the care of 21 sea turtles, initiated new foraging-readiness testing for release candidates in deep-dive tank, and released 14 animals. In 2020, they also initiated a study to better understand how micro-plastics are impacting the sea turtles in their care. In 2016, Florida enacted extensive protection measures. Florida statutes (F.A.C. Rule 68E-1) restrict the take, possession, disturbance, mutilation, destruction, selling, transference, molestation, and harassment of marine turtles, nests or eggs. Protection is also afforded to marine turtle habitat. A specific authorization from commission staff is required to conduct scientific, conservation, or educational activities that directly involve marine turtles in or collected from Florida, their nests, hatchlings or parts thereof, regardless of applicant's possession of any federal permit. In the State of Hawaii, specifically on the Island of Hawai'i (Hawaii County), state representative Faye Hanohano, a Native Hawaiian rights activist, pressed for a measure to delist C. mydas from protected status so that Native Hawaiians could legally harvest the turtles and possibly their eggs as well. The bill, HCR14, was largely overlooked by the media since at that point it was only a local issue. While the bill was passed in the United States House of Representatives, the United States Senate's Committee on Energy and Environment refused to hear it, which meant that the bill did not go on to be heard by the Senate. Country-specific initiatives In addition to management by global entities such as the IUCN and CITES, specific countries around the world have undertaken conservation efforts. The Indonesian island of Bali has traditional uses that were considered sustainable, but have been questioned considering greater demand from the larger and wealthier human population. The harvest was the most intensive in the world. In 1999, Indonesia restricted turtle trade and consumption because of the decreasing population and threat of a tourist boycott. It rejected a request made by Bali Governor I Made Mangku Pastika in November 2009 to set a quota of 1,000 turtles to be killed in Hindu religious ceremonies. While conservationists respect the need for turtles in rituals, they wanted a smaller quota. Multiple protected areas of the Philippines have significant green sea turtle nesting and feeding sites. The most notable is Turtle Islands Wildlife Sanctuary, an UNESCO tentative site which encompasses an entire municipality and one of Southeast Asia's most important green sea turtle nesting sites. Other notable sites include the UNESCO tentative site of El Nido-Taytay Management Resource Protected Area and the UNESCO World Heritage Site of Tubbataha Reefs Natural Park. The species is protected under Republic Act 9147 or the Wildlife Resources Conservation and Protection Act, while the sites where they live and nest are protected under the National Integrated Protected Areas System Act. Ecotourism is one initiative in Sabah, Malaysia. The island of Pulau Selingan is home to a turtle hatchery. Staff people place some of the eggs laid each night in a hatchery to protect them from predators. Incubation takes around sixty days. When the eggs hatch, tourists assist in the release of the baby turtles into the sea. The Hawaiian subpopulation has made a remarkable comeback and is now one focus of ecotourism and has become something of a state mascot. Students of Hawaii Preparatory Academy on the Big Island have tagged thousands of specimens since the early 1990s. In the United Kingdom the species is protected by a Biodiversity Action Plan, due to excess harvesting and marine pollution. The Pakistani-branch of the World Wide Fund for Nature has been initiating projects for secure turtle hatching since the 1980s. However, the population has continued to decline. In the Atlantic, conservation initiatives have centered around Caribbean nesting sites. The Tortuguero nesting beaches in Costa Rica have been the subject of egg-collection limits since the 1950s. The Tortuguero National Park was formally established in 1976, in part, to protect that region's nesting grounds. On Ascension Island, which contains some of the most important nesting beaches, an active conservation program has been implemented. Karumbé has been monitoring foraging and developmental areas of juvenile green turtles in Uruguay since 1999. In Mozambique, there are a number of initiatives to protect sea turtles. In the Primeiras e Segundas, WWF Mozambique has established a turtle tagging and protection program. The archipelago is a vital nesting area for green turtles, including Ilha do Fogo where Fire Island Conservation manage a turtle monitoring programme, and at Celdeira Island, where several nesting females have been tagged. Cayman Turtle Farm located in Grand Cayman in the northwest Caribbean Sea is the first farm to have achieved the second generation of green sea turtles bred, laid, hatched, and raised in captivity. Since its beginning in 1968, the farm has released over 31,000 turtles into the wild, and each year more captive-bred turtles are released into the Caribbean Sea from beaches around the island of Grand Cayman. Captive-bred turtles released from the farm as hatchlings or yearlings with "living tags," have now begun to return to nest on Grand Cayman as adults. On February 19, 2012 the farm released the first 2nd-generation captive-bred green sea turtle equipped with a Position Tracking Transponder, or PTT (also known as a satellite tag). In addition, the farm provides turtle meat products to the local population for whom turtle has been part of the traditional cuisine for centuries. In so doing, the farm curtails the incentive to take turtles from the wild, which over the years in addition to the Cayman Turtle Farm's release of captive-bred turtles has enabled an increase in the number of turtles sighted in the waters around the island of Grand Cayman and nesting on its beaches. In the Pacific, green sea turtles nest on the motu (islets) in the Funafuti Conservation Area, a marine conservation area covering 33 square kilometers (12.74 square miles) of reef, lagoon and motu on the western side of Funafuti atoll in Tuvalu. On Raine Island, up to 100,000 nesting females have been observed in a season, with the cay producing 90% of the region's green turtles. However, the hatching rate declined in the 1990s, and a further decline in the population was threatened by the deaths of thousands of females as they struggled to climb the small sandy cliffs. In addition, as the shape of the island had changed over time, the spread of the beaches outwards had led to greater risk of inundation of the turtle nests. Between 2011 and 2020, a collaborative project by the Queensland Government, BHP (as corporate sponsor), the Great Barrier Reef Marine Park Authority, Great Barrier Reef Foundation, and Wuthathi and Meriam traditional owners, reshaped the island using heavy machinery in a way that gave the female turtles a smoother passage and reduced the risk of nest inundation. A sophisticated monitoring and research system, using 3D modelling, satellite technology and drones was employed, and monitoring continues. , a project called "The Turtle Cooling Project" is being undertaken by scientists from the World Wildlife Fund Australia, University of Queensland, Deakin University and the Queensland Government. It is looking at the effect of global warming on northern green turtle breeding, in particular the effect of producing more male turtles owing to the higher temperatures. They are working in the area around Raine Island, Heron Island and Moulter Cay. Genetics The genome of Chelonia mydas was sequenced in 2013 to examine the development and evolution of the turtle body plan.	Biology and health sciences	Turtles	Animals
5431281	https://en.wikipedia.org/wiki/Circumferentor	Circumferentor	A circumferentor, or surveyor's compass, is an instrument used in surveying to measure horizontal angles. It was superseded by the theodolite in the early 19th century. A circumferentor consists of a circular brass box containing a magnetic needle, which moves freely over a brass circle, or compass divided into 360 degrees. The needle is protected by a glass covering. A pair of sights is located on the North-South axis of the compass. Circumferentors were typically mounted on tripods and rotated on ball-and-socket joints. Circumferentors were made throughout Europe, including in England, France, Italy, and Holland. By the early 19th century, Europeans preferred theodolites to circumferentors. However, the circumferentor remained in common use in mines and in wooded or uncleared areas, such as in America. Usage Measuring angles To measure an angle with a circumferentor, such as angle EKG (Figure 1), place the instrument at K, with the fleur-de-lis in the card towards you. Then direct the sights, until through them you see E; and note the degree pointed at by the south end of the needle, such as 296°. Then, turn the instrument around, with the fleur-de-lis still towards you, and direct the sights to G; note the degree at which the south end of the needle point, such as 182°. Finally, subtract the lesser number, 182, from the greater, 296°; the remainder, 114°, is the number of degrees in the angle EKG. If the remainder is more than 180 degrees, it must be subtracted from 360 degrees. Surveying a region To take the plot of a field, forest, park, etc., with a circumferentor, consider region ABCDEFGHK in Figure 2, an area to be surveyed. Placing the instrument at A, the fleur-de-lis towards you, direct the sights to B; where suppose the south end of the needle cuts 191°; and the ditch, wall, or hedge, measuring with a Gunter's chain, contains 10 chains, 75 links. Placing the instrument at B, direct the sights as before to C; the south end of the needle, e.g. will cut 279°; and the line BC contains 6 chains and 83 links. Then move the instrument to C; turn the sights to measure D, and measure CD as before. In the same manner, proceed to D, E, F, G, H, and lastly to K; still noting the degrees of every bearing, or angle, and the distances of every side. This will result in a table of the following form: From this table, the field is to be plotted, or protracted. Alternative plotting method: An alternative way to plot the area in Figure 2 is to use several angles and only a few measurements and calculate their positions. This could be done by starting at the center point in Figure 2 which is not labeled, but which will be referred to as "Center." Assume each point can be seen from each other point. From the "center" point, sight and record the angle to each point using the sights as described above. Then move to, and measure the distance to, one of the other points referenced, such as point B. At point B, measure the angles to all the other points. Then, move to an additional point such as point F. Again, measure the distance from the center to the point chosen (F). At that point, measure and record the angles to each of the other points as was done at point B. Chose a scale (a ratio between the size of the area to be plotted and the size of the paper on which you will draw the plot) that will allow the plot to fit on your paper and plot the angles and distances. The advantage of this method over the first one above is that there are fewer distance measurements and any errors in angles or distances will not be cumulative; that is, if you use the first survey method, any angle that is slightly off will distort the remainder of the plot. The second method can also be used when it is not possible to measure some of the distances, for example, if there is a water barrier between two of the points. Also, if there are any inaccuracies in the measurements, they will be revealed in the plot because the points plotted from various angles will not coincide. Additional considerations include the number of times the circumferentor must be set up and aligned. With the first method, the instrument must be set up at each point with a compass. With the second method, the initial set up is at "center." After that, for example at point B, the instrument can be set up by aligning the sight with the reciprocal of the angle between "center" and B. Thus, any local change in the magnetic field that would affect the compass would be nullified. Surveyor's double prism A double prism is a device to measure right angles, consisting of two five sided prisms stacked on top of each other and a plumb-bob. It is used to stake out right angles, for example on a construction site.	Technology	Surveying tools	null
2958015	https://en.wikipedia.org/wiki/Philosophy%20of%20artificial%20intelligence	Philosophy of artificial intelligence	The philosophy of artificial intelligence is a branch of the philosophy of mind and the philosophy of computer science that explores artificial intelligence and its implications for knowledge and understanding of intelligence, ethics, consciousness, epistemology, and free will. Furthermore, the technology is concerned with the creation of artificial animals or artificial people (or, at least, artificial creatures; see artificial life) so the discipline is of considerable interest to philosophers. These factors contributed to the emergence of the philosophy of artificial intelligence. The philosophy of artificial intelligence attempts to answer such questions as follows: Can a machine act intelligently? Can it solve any problem that a person would solve by thinking? Are human intelligence and machine intelligence the same? Is the human brain essentially a computer? Can a machine have a mind, mental states, and consciousness in the same sense that a human being can? Can it feel how things are? (i.e. does it have qualia?) Questions like these reflect the divergent interests of AI researchers, cognitive scientists and philosophers respectively. The scientific answers to these questions depend on the definition of "intelligence" and "consciousness" and exactly which "machines" are under discussion. Important propositions in the philosophy of AI include some of the following: Turing's "polite convention": If a machine behaves as intelligently as a human being, then it is as intelligent as a human being. The Dartmouth proposal: "Every aspect of learning or any other feature of intelligence can in principle be so precisely described that a machine can be made to simulate it." Allen Newell and Herbert A. Simon's physical symbol system hypothesis: "A physical symbol system has the necessary and sufficient means of general intelligent action." John Searle's strong AI hypothesis: "The appropriately programmed computer with the right inputs and outputs would thereby have a mind in exactly the same sense human beings have minds." Hobbes' mechanism: "For 'reason' ... is nothing but 'reckoning,' that is adding and subtracting, of the consequences of general names agreed upon for the 'marking' and 'signifying' of our thoughts..." Can a machine display general intelligence? Is it possible to create a machine that can solve all the problems humans solve using their intelligence? This question defines the scope of what machines could do in the future and guides the direction of AI research. It only concerns the behavior of machines and ignores the issues of interest to psychologists, cognitive scientists and philosophers, evoking the question: does it matter whether a machine is really thinking, as a person thinks, rather than just producing outcomes that appear to result from thinking? The basic position of most AI researchers is summed up in this statement, which appeared in the proposal for the Dartmouth workshop of 1956: "Every aspect of learning or any other feature of intelligence can in principle be so precisely described that a machine can be made to simulate it." Arguments against the basic premise must show that building a working AI system is impossible because there is some practical limit to the abilities of computers or that there is some special quality of the human mind that is necessary for intelligent behavior and yet cannot be duplicated by a machine (or by the methods of current AI research). Arguments in favor of the basic premise must show that such a system is possible. It is also possible to sidestep the connection between the two parts of the above proposal. For instance, machine learning, beginning with Turing's infamous child machine proposal, essentially achieves the desired feature of intelligence without a precise design-time description as to how it would exactly work. The account on robot tacit knowledge eliminates the need for a precise description altogether. The first step to answering the question is to clearly define "intelligence". Intelligence Turing test Alan Turing reduced the problem of defining intelligence to a simple question about conversation. He suggests that: if a machine can answer any question posed to it, using the same words that an ordinary person would, then we may call that machine intelligent. A modern version of his experimental design would use an online chat room, where one of the participants is a real person and one of the participants is a computer program. The program passes the test if no one can tell which of the two participants is human. Turing notes that no one (except philosophers) ever asks the question "can people think?" He writes "instead of arguing continually over this point, it is usual to have a polite convention that everyone thinks". Turing's test extends this polite convention to machines: If a machine acts as intelligently as a human being, then it is as intelligent as a human being. One criticism of the Turing test is that it only measures the "humanness" of the machine's behavior, rather than the "intelligence" of the behavior. Since human behavior and intelligent behavior are not exactly the same thing, the test fails to measure intelligence. Stuart J. Russell and Peter Norvig write that "aeronautical engineering texts do not define the goal of their field as 'making machines that fly so exactly like pigeons that they can fool other pigeons'". Intelligence as achieving goals Twenty-first century AI research defines intelligence in terms of goal-directed behavior. It views intelligence as a set of problems that the machine is expected to solve – the more problems it can solve, and the better its solutions are, the more intelligent the program is. AI founder John McCarthy defined intelligence as "the computational part of the ability to achieve goals in the world." Stuart Russell and Peter Norvig formalized this definition using abstract intelligent agents. An "agent" is something which perceives and acts in an environment. A "performance measure" defines what counts as success for the agent. "If an agent acts so as to maximize the expected value of a performance measure based on past experience and knowledge then it is intelligent." Definitions like this one try to capture the essence of intelligence. They have the advantage that, unlike the Turing test, they do not also test for unintelligent human traits such as making typing mistakes. They have the disadvantage that they can fail to differentiate between "things that think" and "things that do not". By this definition, even a thermostat has a rudimentary intelligence. Arguments that a machine can display general intelligence The brain can be simulated Hubert Dreyfus describes this argument as claiming that "if the nervous system obeys the laws of physics and chemistry, which we have every reason to suppose it does, then ... we ... ought to be able to reproduce the behavior of the nervous system with some physical device". This argument, first introduced as early as 1943 and vividly described by Hans Moravec in 1988, is now associated with futurist Ray Kurzweil, who estimates that computer power will be sufficient for a complete brain simulation by the year 2029. A non-real-time simulation of a thalamocortical model that has the size of the human brain (1011 neurons) was performed in 2005, and it took 50 days to simulate 1 second of brain dynamics on a cluster of 27 processors. Even AI's harshest critics (such as Hubert Dreyfus and John Searle) agree that a brain simulation is possible in theory. However, Searle points out that, in principle, anything can be simulated by a computer; thus, bringing the definition to its breaking point leads to the conclusion that any process at all can technically be considered "computation". "What we wanted to know is what distinguishes the mind from thermostats and livers," he writes. Thus, merely simulating the functioning of a living brain would in itself be an admission of ignorance regarding intelligence and the nature of the mind, like trying to build a jet airliner by copying a living bird precisely, feather by feather, with no theoretical understanding of aeronautical engineering. Human thinking is symbol processing In 1963, Allen Newell and Herbert A. Simon proposed that "symbol manipulation" was the essence of both human and machine intelligence. They wrote: "A physical symbol system has the necessary and sufficient means of general intelligent action." This claim is very strong: it implies both that human thinking is a kind of symbol manipulation (because a symbol system is necessary for intelligence) and that machines can be intelligent (because a symbol system is sufficient for intelligence). Another version of this position was described by philosopher Hubert Dreyfus, who called it "the psychological assumption": "The mind can be viewed as a device operating on bits of information according to formal rules." The "symbols" that Newell, Simon and Dreyfus discussed were word-like and high levelsymbols that directly correspond with objects in the world, such as <dog> and <tail>. Most AI programs written between 1956 and 1990 used this kind of symbol. Modern AI, based on statistics and mathematical optimization, does not use the high-level "symbol processing" that Newell and Simon discussed. Arguments against symbol processing These arguments show that human thinking does not consist (solely) of high level symbol manipulation. They do not show that artificial intelligence is impossible, only that more than symbol processing is required. Gödelian anti-mechanist arguments In 1931, Kurt Gödel proved with an incompleteness theorem that it is always possible to construct a "Gödel statement" that a given consistent formal system of logic (such as a high-level symbol manipulation program) could not prove. Despite being a true statement, the constructed Gödel statement is unprovable in the given system. (The truth of the constructed Gödel statement is contingent on the consistency of the given system; applying the same process to a subtly inconsistent system will appear to succeed, but will actually yield a false "Gödel statement" instead.) More speculatively, Gödel conjectured that the human mind can eventually correctly determine the truth or falsity of any well-grounded mathematical statement (including any possible Gödel statement), and that therefore the human mind's power is not reducible to a mechanism. Philosopher John Lucas (since 1961) and Roger Penrose (since 1989) have championed this philosophical anti-mechanist argument. Gödelian anti-mechanist arguments tend to rely on the innocuous-seeming claim that a system of human mathematicians (or some idealization of human mathematicians) is both consistent (completely free of error) and believes fully in its own consistency (and can make all logical inferences that follow from its own consistency, including belief in its Gödel statement) . This is probably impossible for a Turing machine to do (see Halting problem); therefore, the Gödelian concludes that human reasoning is too powerful to be captured by a Turing machine, and by extension, any digital mechanical device. However, the modern consensus in the scientific and mathematical community is that actual human reasoning is inconsistent; that any consistent "idealized version" H of human reasoning would logically be forced to adopt a healthy but counter-intuitive open-minded skepticism about the consistency of H (otherwise H is provably inconsistent); and that Gödel's theorems do not lead to any valid argument that humans have mathematical reasoning capabilities beyond what a machine could ever duplicate. This consensus that Gödelian anti-mechanist arguments are doomed to failure is laid out strongly in Artificial Intelligence: "any attempt to utilize (Gödel's incompleteness results) to attack the computationalist thesis is bound to be illegitimate, since these results are quite consistent with the computationalist thesis." Stuart Russell and Peter Norvig agree that Gödel's argument does not consider the nature of real-world human reasoning. It applies to what can theoretically be proved, given an infinite amount of memory and time. In practice, real machines (including humans) have finite resources and will have difficulty proving many theorems. It is not necessary to be able to prove everything in order to be an intelligent person. Less formally, Douglas Hofstadter, in his Pulitzer Prize winning book Gödel, Escher, Bach: An Eternal Golden Braid, states that these "Gödel-statements" always refer to the system itself, drawing an analogy to the way the Epimenides paradox uses statements that refer to themselves, such as "this statement is false" or "I am lying". But, of course, the Epimenides paradox applies to anything that makes statements, whether it is a machine or a human, even Lucas himself. Consider: Lucas can't assert the truth of this statement. This statement is true but cannot be asserted by Lucas. This shows that Lucas himself is subject to the same limits that he describes for machines, as are all people, and so Lucas's argument is pointless. After concluding that human reasoning is non-computable, Penrose went on to controversially speculate that some kind of hypothetical non-computable processes involving the collapse of quantum mechanical states give humans a special advantage over existing computers. Existing quantum computers are only capable of reducing the complexity of Turing computable tasks and are still restricted to tasks within the scope of Turing machines. . By Penrose and Lucas's arguments, the fact that quantum computers are only able to complete Turing computable tasks implies that they cannot be sufficient for emulating the human mind. Therefore, Penrose seeks for some other process involving new physics, for instance quantum gravity which might manifest new physics at the scale of the Planck mass via spontaneous quantum collapse of the wave function. These states, he suggested, occur both within neurons and also spanning more than one neuron. However, other scientists point out that there is no plausible organic mechanism in the brain for harnessing any sort of quantum computation, and furthermore that the timescale of quantum decoherence seems too fast to influence neuron firing. Dreyfus: the primacy of implicit skills Hubert Dreyfus argued that human intelligence and expertise depended primarily on fast intuitive judgements rather than step-by-step symbolic manipulation, and argued that these skills would never be captured in formal rules. Dreyfus's argument had been anticipated by Turing in his 1950 paper Computing machinery and intelligence, where he had classified this as the "argument from the informality of behavior." Turing argued in response that, just because we do not know the rules that govern a complex behavior, this does not mean that no such rules exist. He wrote: "we cannot so easily convince ourselves of the absence of complete laws of behaviour ... The only way we know of for finding such laws is scientific observation, and we certainly know of no circumstances under which we could say, 'We have searched enough. There are no such laws.'" Russell and Norvig point out that, in the years since Dreyfus published his critique, progress has been made towards discovering the "rules" that govern unconscious reasoning. The situated movement in robotics research attempts to capture our unconscious skills at perception and attention. Computational intelligence paradigms, such as neural nets, evolutionary algorithms and so on are mostly directed at simulated unconscious reasoning and learning. Statistical approaches to AI can make predictions which approach the accuracy of human intuitive guesses. Research into commonsense knowledge has focused on reproducing the "background" or context of knowledge. In fact, AI research in general has moved away from high level symbol manipulation, towards new models that are intended to capture more of our intuitive reasoning. Cognitive science and psychology eventually came to agree with Dreyfus' description of human expertise. Daniel Kahnemann and others developed a similar theory where they identified two "systems" that humans use to solve problems, which he called "System 1" (fast intuitive judgements) and "System 2" (slow deliberate step by step thinking). Although Dreyfus' views have been vindicated in many ways, the work in cognitive science and in AI was in response to specific problems in those fields and was not directly influenced by Dreyfus. Historian and AI researcher Daniel Crevier wrote that "time has proven the accuracy and perceptiveness of some of Dreyfus's comments. Had he formulated them less aggressively, constructive actions they suggested might have been taken much earlier." Can a machine have a mind, consciousness, and mental states? This is a philosophical question, related to the problem of other minds and the hard problem of consciousness. The question revolves around a position defined by John Searle as "strong AI": A physical symbol system can have a mind and mental states. Searle distinguished this position from what he called "weak AI": A physical symbol system can act intelligently. Searle introduced the terms to isolate strong AI from weak AI so he could focus on what he thought was the more interesting and debatable issue. He argued that even if we assume that we had a computer program that acted exactly like a human mind, there would still be a difficult philosophical question that needed to be answered. Neither of Searle's two positions are of great concern to AI research, since they do not directly answer the question "can a machine display general intelligence?" (unless it can also be shown that consciousness is necessary for intelligence). Turing wrote "I do not wish to give the impression that I think there is no mystery about consciousness… [b]ut I do not think these mysteries necessarily need to be solved before we can answer the question [of whether machines can think]." Russell and Norvig agree: "Most AI researchers take the weak AI hypothesis for granted, and don't care about the strong AI hypothesis." There are a few researchers who believe that consciousness is an essential element in intelligence, such as Igor Aleksander, Stan Franklin, Ron Sun, and Pentti Haikonen, although their definition of "consciousness" strays very close to "intelligence". (See artificial consciousness.) Before we can answer this question, we must be clear what we mean by "minds", "mental states" and "consciousness". Consciousness, minds, mental states, meaning The words "mind" and "consciousness" are used by different communities in different ways. Some new age thinkers, for example, use the word "consciousness" to describe something similar to Bergson's "élan vital": an invisible, energetic fluid that permeates life and especially the mind. Science fiction writers use the word to describe some essential property that makes us human: a machine or alien that is "conscious" will be presented as a fully human character, with intelligence, desires, will, insight, pride and so on. (Science fiction writers also use the words "sentience", "sapience", "self-awareness" or "ghost"—as in the Ghost in the Shell manga and anime series—to describe this essential human property). For others , the words "mind" or "consciousness" are used as a kind of secular synonym for the soul. For philosophers, neuroscientists and cognitive scientists, the words are used in a way that is both more precise and more mundane: they refer to the familiar, everyday experience of having a "thought in your head", like a perception, a dream, an intention or a plan, and to the way we see something, know something, mean something or understand something. "It's not hard to give a commonsense definition of consciousness" observes philosopher John Searle. What is mysterious and fascinating is not so much what it is but how it is: how does a lump of fatty tissue and electricity give rise to this (familiar) experience of perceiving, meaning or thinking? Philosophers call this the hard problem of consciousness. It is the latest version of a classic problem in the philosophy of mind called the "mind-body problem". A related problem is the problem of meaning or understanding (which philosophers call "intentionality"): what is the connection between our thoughts and what we are thinking about (i.e. objects and situations out in the world)? A third issue is the problem of experience (or "phenomenology"): If two people see the same thing, do they have the same experience? Or are there things "inside their head" (called "qualia") that can be different from person to person? Neurobiologists believe all these problems will be solved as we begin to identify the neural correlates of consciousness: the actual relationship between the machinery in our heads and its collective properties; such as the mind, experience and understanding. Some of the harshest critics of artificial intelligence agree that the brain is just a machine, and that consciousness and intelligence are the result of physical processes in the brain. The difficult philosophical question is this: can a computer program, running on a digital machine that shuffles the binary digits of zero and one, duplicate the ability of the neurons to create minds, with mental states (like understanding or perceiving), and ultimately, the experience of consciousness? Arguments that a computer cannot have a mind and mental states Searle's Chinese room John Searle asks us to consider a thought experiment: suppose we have written a computer program that passes the Turing test and demonstrates general intelligent action. Suppose, specifically that the program can converse in fluent Chinese. Write the program on 3x5 cards and give them to an ordinary person who does not speak Chinese. Lock the person into a room and have him follow the instructions on the cards. He will copy out Chinese characters and pass them in and out of the room through a slot. From the outside, it will appear that the Chinese room contains a fully intelligent person who speaks Chinese. The question is this: is there anyone (or anything) in the room that understands Chinese? That is, is there anything that has the mental state of understanding, or which has conscious awareness of what is being discussed in Chinese? The man is clearly not aware. The room cannot be aware. The cards certainly are not aware. Searle concludes that the Chinese room, or any other physical symbol system, cannot have a mind. Searle goes on to argue that actual mental states and consciousness require (yet to be described) "actual physical-chemical properties of actual human brains." He argues there are special "causal properties" of brains and neurons that gives rise to minds: in his words "brains cause minds." Related arguments: Leibniz' mill, Davis's telephone exchange, Block's Chinese nation and Blockhead Gottfried Leibniz made essentially the same argument as Searle in 1714, using the thought experiment of expanding the brain until it was the size of a mill. In 1974, Lawrence Davis imagined duplicating the brain using telephone lines and offices staffed by people, and in 1978 Ned Block envisioned the entire population of China involved in such a brain simulation. This thought experiment is called "the Chinese Nation" or "the Chinese Gym". Ned Block also proposed his Blockhead argument, which is a version of the Chinese room in which the program has been re-factored into a simple set of rules of the form "see this, do that", removing all mystery from the program. Responses to the Chinese room Responses to the Chinese room emphasize several different points. The systems reply and the virtual mind reply: This reply argues that the system, including the man, the program, the room, and the cards, is what understands Chinese. Searle claims that the man in the room is the only thing which could possibly "have a mind" or "understand", but others disagree, arguing that it is possible for there to be two minds in the same physical place, similar to the way a computer can simultaneously "be" two machines at once: one physical (like a Macintosh) and one "virtual" (like a word processor). Speed, power and complexity replies: Several critics point out that the man in the room would probably take millions of years to respond to a simple question, and would require "filing cabinets" of astronomical proportions. This brings the clarity of Searle's intuition into doubt. Robot reply: To truly understand, some believe the Chinese Room needs eyes and hands. Hans Moravec writes: "If we could graft a robot to a reasoning program, we wouldn't need a person to provide the meaning anymore: it would come from the physical world." Brain simulator reply: What if the program simulates the sequence of nerve firings at the synapses of an actual brain of an actual Chinese speaker? The man in the room would be simulating an actual brain. This is a variation on the "systems reply" that appears more plausible because "the system" now clearly operates like a human brain, which strengthens the intuition that there is something besides the man in the room that could understand Chinese. Other minds reply and the epiphenomena reply: Several people have noted that Searle's argument is just a version of the problem of other minds, applied to machines. Since it is difficult to decide if people are "actually" thinking, we should not be surprised that it is difficult to answer the same question about machines. A related question is whether "consciousness" (as Searle understands it) exists. Searle argues that the experience of consciousness cannot be detected by examining the behavior of a machine, a human being or any other animal. Daniel Dennett points out that natural selection cannot preserve a feature of an animal that has no effect on the behavior of the animal, and thus consciousness (as Searle understands it) cannot be produced by natural selection. Therefore, either natural selection did not produce consciousness, or "strong AI" is correct in that consciousness can be detected by suitably designed Turing test. Is thinking a kind of computation? The computational theory of mind or "computationalism" claims that the relationship between mind and brain is similar (if not identical) to the relationship between a running program (software) and a computer (hardware). The idea has philosophical roots in Hobbes (who claimed reasoning was "nothing more than reckoning"), Leibniz (who attempted to create a logical calculus of all human ideas), Hume (who thought perception could be reduced to "atomic impressions") and even Kant (who analyzed all experience as controlled by formal rules). The latest version is associated with philosophers Hilary Putnam and Jerry Fodor. This question bears on our earlier questions: if the human brain is a kind of computer then computers can be both intelligent and conscious, answering both the practical and philosophical questions of AI. In terms of the practical question of AI ("Can a machine display general intelligence?"), some versions of computationalism make the claim that (as Hobbes wrote): Reasoning is nothing but reckoning. In other words, our intelligence derives from a form of calculation, similar to arithmetic. This is the physical symbol system hypothesis discussed above, and it implies that artificial intelligence is possible. In terms of the philosophical question of AI ("Can a machine have mind, mental states and consciousness?"), most versions of computationalism claim that (as Stevan Harnad characterizes it): Mental states are just implementations of (the right) computer programs. This is John Searle's "strong AI" discussed above, and it is the real target of the Chinese room argument (according to Harnad). Other related questions Can a machine have emotions? If "emotions" are defined only in terms of their effect on behavior or on how they function inside an organism, then emotions can be viewed as a mechanism that an intelligent agent uses to maximize the utility of its actions. Given this definition of emotion, Hans Moravec believes that "robots in general will be quite emotional about being nice people". Fear is a source of urgency. Empathy is a necessary component of good human computer interaction. He says robots "will try to please you in an apparently selfless manner because it will get a thrill out of this positive reinforcement. You can interpret this as a kind of love." Daniel Crevier writes "Moravec's point is that emotions are just devices for channeling behavior in a direction beneficial to the survival of one's species." Can a machine be self-aware? "Self-awareness", as noted above, is sometimes used by science fiction writers as a name for the essential human property that makes a character fully human. Turing strips away all other properties of human beings and reduces the question to "can a machine be the subject of its own thought?" Can it think about itself? Viewed in this way, a program can be written that can report on its own internal states, such as a debugger. Can a machine be original or creative? Turing reduces this to the question of whether a machine can "take us by surprise" and argues that this is obviously true, as any programmer can attest. He notes that, with enough storage capacity, a computer can behave in an astronomical number of different ways. It must be possible, even trivial, for a computer that can represent ideas to combine them in new ways. (Douglas Lenat's Automated Mathematician, as one example, combined ideas to discover new mathematical truths.) Kaplan and Haenlein suggest that machines can display scientific creativity, while it seems likely that humans will have the upper hand where artistic creativity is concerned. In 2009, scientists at Aberystwyth University in Wales and the U.K's University of Cambridge designed a robot called Adam that they believe to be the first machine to independently come up with new scientific findings. Also in 2009, researchers at Cornell developed Eureqa, a computer program that extrapolates formulas to fit the data inputted, such as finding the laws of motion from a pendulum's motion. Can a machine be benevolent or hostile? This question (like many others in the philosophy of artificial intelligence) can be presented in two forms. "Hostility" can be defined in terms function or behavior, in which case "hostile" becomes synonymous with "dangerous". Or it can be defined in terms of intent: can a machine "deliberately" set out to do harm? The latter is the question "can a machine have conscious states?" (such as intentions) in another form. The question of whether highly intelligent and completely autonomous machines would be dangerous has been examined in detail by futurists (such as the Machine Intelligence Research Institute). The obvious element of drama has also made the subject popular in science fiction, which has considered many differently possible scenarios where intelligent machines pose a threat to mankind; see Artificial intelligence in fiction. One issue is that machines may acquire the autonomy and intelligence required to be dangerous very quickly. Vernor Vinge has suggested that over just a few years, computers will suddenly become thousands or millions of times more intelligent than humans. He calls this "the Singularity". He suggests that it may be somewhat or possibly very dangerous for humans. This is discussed by a philosophy called Singularitarianism. In 2009, academics and technical experts attended a conference to discuss the potential impact of robots and computers and the impact of the hypothetical possibility that they could become self-sufficient and able to make their own decisions. They discussed the possibility and the extent to which computers and robots might be able to acquire any level of autonomy, and to what degree they could use such abilities to possibly pose any threat or hazard. They noted that some machines have acquired various forms of semi-autonomy, including being able to find power sources on their own and being able to independently choose targets to attack with weapons. They also noted that some computer viruses can evade elimination and have achieved "cockroach intelligence". They noted that self-awareness as depicted in science-fiction is probably unlikely, but that there were other potential hazards and pitfalls. Some experts and academics have questioned the use of robots for military combat, especially when such robots are given some degree of autonomous functions. The US Navy has funded a report which indicates that as military robots become more complex, there should be greater attention to implications of their ability to make autonomous decisions. The President of the Association for the Advancement of Artificial Intelligence has commissioned a study to look at this issue. They point to programs like the Language Acquisition Device which can emulate human interaction. Some have suggested a need to build "Friendly AI", a term coined by Eliezer Yudkowsky, meaning that the advances which are already occurring with AI should also include an effort to make AI intrinsically friendly and humane. Can a machine imitate all human characteristics? Turing said "It is customary ... to offer a grain of comfort, in the form of a statement that some peculiarly human characteristic could never be imitated by a machine. ... I cannot offer any such comfort, for I believe that no such bounds can be set." Turing noted that there are many arguments of the form "a machine will never do X", where X can be many things, such as: Be kind, resourceful, beautiful, friendly, have initiative, have a sense of humor, tell right from wrong, make mistakes, fall in love, enjoy strawberries and cream, make someone fall in love with it, learn from experience, use words properly, be the subject of its own thought, have as much diversity of behaviour as a man, do something really new. Turing argues that these objections are often based on naive assumptions about the versatility of machines or are "disguised forms of the argument from consciousness". Writing a program that exhibits one of these behaviors "will not make much of an impression." All of these arguments are tangential to the basic premise of AI, unless it can be shown that one of these traits is essential for general intelligence. Can a machine have a soul? Finally, those who believe in the existence of a soul may argue that "Thinking is a function of man's immortal soul." Alan Turing called this "the theological objection". He writes: In attempting to construct such machines we should not be irreverently usurping His power of creating souls, any more than we are in the procreation of children: rather we are, in either case, instruments of His will providing mansions for the souls that He creates.The discussion on the topic has been reignited as a result of recent claims made by Google's LaMDA artificial intelligence system that it is sentient and had a "soul". LaMDA (Language Model for Dialogue Applications) is an artificial intelligence system that creates chatbots—AI robots designed to communicate with humans—by gathering vast amounts of text from the internet and using algorithms to respond to queries in the most fluid and natural way possible. The transcripts of conversations between scientists and LaMDA reveal that the AI system excels at this, providing answers to challenging topics about the nature of emotions, generating Aesop-style fables on the moment, and even describing its alleged fears. Pretty much all philosophers doubt LaMDA's sentience. Views on the role of philosophy Some scholars argue that the AI community's dismissal of philosophy is detrimental. In the Stanford Encyclopedia of Philosophy, some philosophers argue that the role of philosophy in AI is underappreciated. Physicist David Deutsch argues that without an understanding of philosophy or its concepts, AI development would suffer from a lack of progress. Conferences and literature The main conference series on the issue is "Philosophy and Theory of AI" (PT-AI), run by Vincent C. Müller. The main bibliography on the subject, with several sub-sections, is on PhilPapers. A recent survey for Philosophy of AI is Müller (2023).	Technology	Artificial intelligence concepts	null
2959226	https://en.wikipedia.org/wiki/Acacia	Acacia	Acacia, commonly known as wattles or acacias, is a genus of about of shrubs and trees in the subfamily Mimosoideae of the pea family Fabaceae. Initially, it comprised a group of plant species native to Africa, South America, and Australasia, but is now reserved for species mainly from Australia, with others from New Guinea, Southeast Asia, and the Indian Ocean. The genus name is Neo-Latin, borrowed from the Greek (), a term used in antiquity to describe a preparation extracted from Vachellia nilotica, the original type species. A number of species of Acacia have been introduced to various parts of the world, and two million hectares of commercial plantations have been established. Description Plants in the genus Acacia are shrubs or trees with bipinnate leaves, the mature leaves sometimes reduced to phyllodes or rarely absent. There are 2 small stipules at the base of the leaf, but sometimes fall off as the leaf matures. The flowers are borne in spikes or cylindrical heads, sometimes singly, in pairs or in racemes in the axils of leaves or phyllodes, sometimes in panicles on the ends of branches. Each spike or cylindrical head has many small golden-yellow to pale creamy-white flowers, each with 4 or 5 sepals and petals, more than 10 stamens, and a thread-like style that is longer than the stamens. The fruit is a variably-shaped pod, sometimes flat or cylindrical, containing seeds with a fleshy aril on the end. Taxonomy The genus was first validly named in 1754 by Philip Miller in The Gardeners Dictionary. In 1913 Nathaniel Lord Britton and Addison Brown selected Mimosa scorpioides (≡ Acacia scorpioides () = Acacia nilotica () ), a species from Africa, as the lectotype of the name. Etymology The genus name comes from Neo-Latin; Gaspard Bauhin in his book Pinax (1623) writes it coming from Pedanius Dioscorides who uses the name akakia for species Vachellia nilotica, the original type species growing in Egypt, from akakis meaning "point". The origin of "wattle" may be an Old Teutonic word meaning "to weave". From around 700 AD, was used in Old English to refer to the flexible woody vines, branches, and sticks which were interwoven to form walls, roofs, and fences. Since about 1810 it has been used as the common name for the Australian legume trees and shrubs that can provide these branches. History Genus Acacia was considered to contain some leading to 1986. That year, Leslie Pedley questioned the monophyletic nature of the genus, and proposed a split into three genera: Acacia sensu stricto (161 species), Senegalia (231 species) and Racosperma (960 species), the last name first proposed in 1829 by Carl Friedrich Philipp von Martius as the name of a section in Acacia, but raised to generic rank in 1835. In 2003, Pedley published a paper with 834 new combinations in Racosperma for species, most of which were formerly placed in Acacia. All but 10 of these species are native to Australasia, where it constitutes the largest plant genus. In the early 2000s, it had become evident that the genus as it stood was not monophyletic and that several divergent lineages needed to be placed in separate genera. It turned out that one lineage comprising over 900 species mainly native to Australia, New Guinea, and Indonesia was not closely related to the much smaller group of African lineage that contained A. nilotica – the type species. This meant that the Australasian lineage (by far the most prolific in number of species) would need to be renamed. Pedley's proposed name of Racosperma for this group had received little acclaim in the botanical community. Australian botanists proposed a less disruptive solution, setting a different type species for Acacia (A. penninervis) and allowing this largest number of species to remain in Acacia, resulting in the two pan-tropical lineages being renamed Vachellia and Senegalia, and the two endemic American lineages renamed Acaciella and Mariosousa. In 2003, Anthony Orchard and Bruce Maslin filed a proposal to conserve the name Acacia with a different type, in order to retain the Australasian group of species in the genus Acacia. Following a controversial decision to choose a new type for Acacia in 2005, the Australian component of Acacia s.l. now retains the name Acacia. At the 2011 International Botanical Congress held in Melbourne, Australia, the decision to use the name Acacia, rather than the proposed Racosperma for this genus, was upheld. Other Acacia s.l. taxa continue to be called Acacia by those who choose to consider the entire group as one genus. The Australian species of the genus Paraserianthes s.l. (namely P. lophantha) are deemed its closest relatives. The nearest relatives of Acacia and Paraserianthes s.l. in turn include the Australian and South East Asian genera Archidendron, Archidendropsis, Pararchidendron and Wallaceodendron, all of the tribe Ingeae. Species The names of more than 1,080 species of Acacia, mostly native to Australia, have been accepted by Plants of the World Online as at January 2025. Fossil record An Acacia-like long fossil seed pod has been described from the Eocene of the Paris Basin. Acacia-like fossil pods under the name Leguminocarpon are known from late Oligocene deposits at different sites in Hungary. Seed pod fossils of †Acacia parschlugiana and †Acacia cyclosperma are known from Tertiary deposits in Switzerland. †Acacia colchica has been described from the Miocene of West Georgia. Pliocene fossil pollen of an Acacia sp. has been described from West Georgia (including Abkhazia). The oldest fossil Acacia pollen in Australia are recorded as being from the late Oligocene epoch, 25 million years ago. Distribution and habitat Species of Acacia occurs in all Australian states and territories, and on its nearby islands. About 20 species occur naturally outside Australia and also occur in Australia. One species (Acacia koa) is native to Hawaii and one (Acacia heterophylla) is native to Mauritius and Réunion in the Indian Ocean. They are present in all terrestrial habitats, including alpine settings, rainforests, woodlands, grasslands, coastal dunes and deserts. In drier woodlands or forests they are an important component of the understory. Elsewhere they may be dominant, as in the Brigalow Belt, Myall woodlands and the eremaean Mulga woodlands. In Australia, Acacia forest is the second most common forest type after eucalypt forest, covering or 8% of total forest area. Acacia is also the nation's largest genus of flowering plants with almost found. Ecology Acacia is a common food source and host plant for butterflies of the genus Jalmenus. The imperial hairstreak, Jalmenus evagoras, feeds on at least 25 acacia species. Many reptiles feed on the sap, such as the native house gecko in Australia. The sap is also consumed by bugs (Hemiptera), such as Hackerobrachys viridiventris and Sextius virescens. Toxicity Some species of acacia contain psychoactive alkaloids, and some contain potassium fluoroacetate. Uses The seed pods, flowers, and young leaves are generally edible either raw or cooked. Aboriginal Australians have traditionally harvested the seeds of some species, to be ground into flour and eaten as a paste or baked into a cake. Wattleseeds contain as much as 25% more protein than common cereals, and they store well for long periods due to the hard seed coats. In addition to consuming the edible seed and gum, Aboriginal people also employed the timber for implements, weapons, fuel and musical instruments. A number of species, most notably Acacia mangium (hickory wattle), A. mearnsii (black wattle) and A. saligna (coojong), are economically important and are widely planted globally for wood products, tannin, firewood and fodder. A. melanoxylon (blackwood) and A. aneura (mulga) supply some of the most attractive timbers in the genus. Black wattle bark supported the tanning industries of several countries, and may supply tannins for production of waterproof adhesives. In Vietnam, Acacia is used in plantations of non-native species that are regularly clear-cut for paper or timber uses. Wattle bark collected in Australia in the 19th century was exported to Europe where it was used in the tanning process. One ton of wattle or mimosa bark contained about of pure tannin. The gum of some species may be used as a substitute for gum arabic, known as Australian gum or wattle gum. Cultivation Some species of acacia – notably Acacia baileyana, A. dealbata and A. pravissima – are cultivated as ornamental garden plants. The 1889 publication Useful Native Plants of Australia describes various uses for eating.	Biology and health sciences	Fabales	Plants
2959657	https://en.wikipedia.org/wiki/Epidermis%20%28botany%29	Epidermis (botany)	The epidermis (from the Greek ἐπιδερμίς, meaning "over-skin") is a single layer of cells that covers the leaves, flowers, roots and stems of plants. It forms a boundary between the plant and the external environment. The epidermis serves several functions: it protects against water loss, regulates gas exchange, secretes metabolic compounds, and (especially in roots) absorbs water and mineral nutrients. The epidermis of most leaves shows dorsoventral anatomy: the upper (adaxial) and lower (abaxial) surfaces have somewhat different construction and may serve different functions. Woody stems and some other stem structures such as potato tubers produce a secondary covering called the periderm that replaces the epidermis as the protective covering. Description The epidermis is the outermost cell layer of the primary plant body. In some older works the cells of the leaf epidermis have been regarded as specialized parenchyma cells, but the established modern preference has long been to classify the epidermis as dermal tissue, whereas parenchyma is classified as ground tissue. The epidermis is the main component of the dermal tissue system of leaves (diagrammed below), and also stems, roots, flowers, fruits, and seeds; it is usually transparent (epidermal cells have fewer chloroplasts or lack them completely, except for the guard cells.) The cells of the epidermis are structurally and functionally variable. Most plants have an epidermis that is a single cell layer thick. Some plants like Ficus elastica and Peperomia, which have a periclinal cellular division within the protoderm of the leaves, have an epidermis with multiple cell layers. Epidermal cells are tightly linked to each other and provide mechanical strength and protection to the plant. Particularly, wavy pavement cells are suggested to play a pivotal role in preventing or guiding cracks in the epidermis. The walls of the epidermal cells of the above-ground parts of plants contain cutin, and are covered with a cuticle. The cuticle reduces water loss to the atmosphere, it is sometimes covered with wax in smooth sheets, granules, plates, tubes, or filaments. The wax layers give some plants a whitish or bluish surface color. Surface wax acts as a moisture barrier and protects the plant from intense sunlight and wind. The epidermal tissue includes several differentiated cell types: epidermal cells, guard cells, subsidiary cells, and epidermal hairs (trichomes). The epidermal cells are the most numerous, largest, and least specialized. These are typically more elongated in the leaves of monocots than in those of dicots. Trichomes or hairs grow out from the epidermis in many species. In the root epidermis, epidermal hairs termed root hairs are common and are specialized for the absorption of water and mineral nutrients. In plants with secondary growth, the epidermis of roots and stems is usually replaced by a periderm through the action of a cork cambium. Stoma complex The leaf and stem epidermis is covered with pores called stomata (sing; stoma), part of a stoma complex consisting of a pore surrounded on each side by chloroplast-containing guard cells, and two to four subsidiary cells that lack chloroplasts. The stomata complex regulates the exchange of gases and water vapor between the outside air and the interior of the leaf. Typically, the stomata are more numerous over the abaxial (lower) epidermis of the leaf than the (adaxial) upper epidermis. An exception is floating leaves where most or all stomata are on the upper surface. Vertical leaves, such as those of many grasses, often have roughly equal numbers of stomata on both surfaces. The stoma is bounded by two guard cells. The guard cells differ from the epidermal cells in the following aspects: The guard cells are bean-shaped in surface view, while the epidermal cells are irregular in shape The guard cells contain chloroplasts, so they can manufacture food by photosynthesis (The epidermal cells of terrestrial plants do not contain chloroplasts) Guard cells are the only epidermal cells that can make sugar. According to one theory, in sunlight, the concentration of potassium ions (K+) increases in the guard cells. This, together with the sugars formed, lowers the water potential in the guard cells. As a result, water from other cells enters the guard cells by osmosis so they swell and become turgid. Because the guard cells have a thicker cellulose wall on one side of the cell, i.e. the side around the stomatal pore, the swollen guard cells become curved and pull the stomata open. At night, the sugar is used up and water leaves the guard cells, so they become flaccid and the stomatal pore closes. In this way, they reduce the amount of water vapor escaping from the leaf. Cell differentiation in the epidermis The plant epidermis consists of three main cell types: pavement cells, guard cells and their subsidiary cells that surround the stomata and trichomes, otherwise known as leaf hairs. The epidermis of petals also form a variation of trichomes called conical cells. Trichomes develop at a distinct phase during leaf development, under the control of two major trichome specification genes: TTG and GL1. The process may be controlled by the plant hormones gibberellins, and even if not completely controlled, gibberellins certainly have an effect on the development of the leaf hairs. GL1 causes endoreplication, the replication of DNA without subsequent cell division as well as cell expansion. GL1 turns on the expression of a second gene for trichome formation, GL2, which controls the final stages of trichome formation causing the cellular outgrowth. Arabidopsis thaliana uses the products of inhibitory genes to control the patterning of trichomes, such as TTG and TRY. The products of these genes will diffuse into the lateral cells, preventing them from forming trichomes and in the case of TRY promoting the formation of pavement cells. Expression of the gene MIXTA, or its analogue in other species, later in the process of cellular differentiation will cause the formation of conical cells over trichomes. MIXTA is a transcription factor. Stomatal patterning is a much more controlled process, as the stoma affects the plant's water retention and respiration capabilities. As a consequence of these important functions, differentiation of cells to form stomata is also subject to environmental conditions to a much greater degree than other epidermal cell types. Stomata are pores in the plant epidermis that are surrounded by two guard cells, which control the opening and closing of the aperture. These guard cells are in turn surrounded by subsidiary cells which provide a supporting role for the guard cells. Stomata begin as stomatal meristemoids. The process differs between dicots and monocots. Spacing is thought to be essentially random in dicots though mutants do show it is under some form of genetic control, but it is more controlled in monocots, where stomata arise from specific asymmetric divisions of protoderm cells. The smaller of the two cells produced becomes the guard mother cells. Adjacent epidermal cells will also divide asymmetrically to form the subsidiary cells. Because stomata play such an important role in the plants' survival, collecting information on their differentiation is difficult by the traditional means of genetic manipulation, as stomatal mutants tend to be unable to survive. Thus the control of the process is not well understood. Some genes have been identified. TMM is thought to control the timing of stomatal initiation specification and FLP is thought to be involved in preventing the further division of the guard cells once they are formed. Environmental conditions affect the development of stomata, in particular, their density on the leaf surface. It is thought that plant hormones, such as ethylene and cytokines, control the stomatal developmental response to the environmental conditions. Accumulation of these hormones appears to cause increased stomatal density such as when the plants are kept in closed environments.	Biology and health sciences	Plant tissues	null
314402	https://en.wikipedia.org/wiki/Liquid%20oxygen	Liquid oxygen	Liquid oxygen, sometimes abbreviated as LOX or LOXygen, is a clear cyan liquid form of dioxygen . It was used as the oxidizer in the first liquid-fueled rocket invented in 1926 by Robert H. Goddard, an application which is ongoing. Physical properties Liquid oxygen has a clear cyan color and is strongly paramagnetic: it can be suspended between the poles of a powerful horseshoe magnet. Liquid oxygen has a density of , slightly denser than liquid water, and is cryogenic with a freezing point of and a boiling point of at . Liquid oxygen has an expansion ratio of 1:861 and because of this, it is used in some commercial and military aircraft as a transportable source of breathing oxygen. Because of its cryogenic nature, liquid oxygen can cause the materials it touches to become extremely brittle. Liquid oxygen is also a very powerful oxidizing agent: organic materials will burn rapidly and energetically in liquid oxygen. Further, if soaked in liquid oxygen, some materials such as coal briquettes, carbon black, etc., can detonate unpredictably from sources of ignition such as flames, sparks or impact from light blows. Petrochemicals, including asphalt, often exhibit this behavior. The tetraoxygen molecule (O4) was predicted in 1924 by Gilbert N. Lewis, who proposed it to explain why liquid oxygen defied Curie's law. Modern computer simulations indicate that, although there are no stable O4 molecules in liquid oxygen, O2 molecules do tend to associate in pairs with antiparallel spins, forming transient O4 units. Liquid nitrogen has a lower boiling point at −196 °C (77 K) than oxygen's −183 °C (90 K), and vessels containing liquid nitrogen can condense oxygen from air: when most of the nitrogen has evaporated from such a vessel, there is a risk that liquid oxygen remaining can react violently with organic material. Conversely, liquid nitrogen or liquid air can be oxygen-enriched by letting it stand in open air; atmospheric oxygen dissolves in it, while nitrogen evaporates preferentially. The surface tension of liquid oxygen at its normal pressure boiling point is . Uses In commerce, liquid oxygen is classified as an industrial gas and is widely used for industrial and medical purposes. Liquid oxygen is obtained from the oxygen found naturally in air by fractional distillation in a cryogenic air separation plant. Air forces have long recognized the strategic importance of liquid oxygen, both as an oxidizer and as a supply of gaseous oxygen for breathing in hospitals and high-altitude aircraft flights. In 1985, the USAF started a program of building its own oxygen-generation facilities at all major consumption bases. In rocket propellant Liquid oxygen is the most common cryogenic liquid oxidizer propellant for spacecraft rocket applications, usually in combination with liquid hydrogen, kerosene or methane. Liquid oxygen was used in the first liquid fueled rocket. The World War II V-2 missile also used liquid oxygen under the name A-Stoff and Sauerstoff. In the 1950s, during the Cold War both the United States' Redstone and Atlas rockets, and the Soviet R-7 Semyorka used liquid oxygen. Later, in the 1960s and 1970s, the ascent stages of the Apollo Saturn rockets, and the Space Shuttle main engines used liquid oxygen. As of 2024, many active rockets use liquid oxygen: Chinese space program CASC: Long March 5, Long March 6, Long March 7, Long March 8, Long March 12, Long March 9 (under development), Long March 10 (under development) Galactic Energy: Pallas-1 (under development) i-Space: Hyperbola-3 (under development) LandSpace: Zhuque-2 Orienspace: Gravity-2 (under development) Space Pioneer: Tianlong-2 European Space Agency: Ariane 6 Indian Space Research Organisation: GSLV JAXA (Japan): H-IIA, H3 Korea Aerospace Research Institute: Naro-1, Nuri Roscosmos (Russia): Soyuz-2, Angara United States Blue Origin: New Shepard, New Glenn (under development) Firefly Aerospace: Firefly Alpha NASA: Space Launch System Northrop Grumman: Antares 300 (under development) Rocket Lab: Electron, Neutron (under development) SpaceX: Falcon 9, Falcon Heavy, Starship United Launch Alliance: Atlas V, Vulcan History By 1845, Michael Faraday had managed to liquefy most gases then known to exist. Six gases, however, resisted every attempt at liquefaction and were known at the time as "permanent gases". They were oxygen, hydrogen, nitrogen, carbon monoxide, methane, and nitric oxide. In 1877, Louis Paul Cailletet in France and Raoul Pictet in Switzerland succeeded in producing the first droplets of liquid air. In 1883, Polish professors Zygmunt Wróblewski and Karol Olszewski produced the first measurable quantity of liquid oxygen.	Physical sciences	Group 16	Chemistry
314494	https://en.wikipedia.org/wiki/Libration	Libration	In lunar astronomy, libration is the cyclic variation in the apparent position of the Moon that is perceived by observers on the Earth observers and caused by changes between the orbital and rotational planes of the moon. It causes an observer to see slightly different hemispheres of the surface at different times. It is similar in both cause and effect to the changes in the Moon's apparent size because of changes in distance. It is caused by three mechanisms detailed below, two of which cause a relatively tiny physical libration via tidal forces exerted by the Earth. Such true librations are known as well for other moons with locked rotation. The quite different phenomenon of a trojan asteroid's movement has been called Trojan libration, and Trojan libration point means Lagrangian point. Lunar libration The Moon keeps one hemisphere of itself facing the Earth because of tidal locking. Therefore, the first view of the far side of the Moon was not possible until the Soviet probe Luna 3 reached the Moon on October 7, 1959, and further lunar exploration by the United States and the Soviet Union. This simple picture is only approximately true since over time, slightly more than half (about 59% in total) of the Moon's surface is seen from Earth because of libration. Lunar libration arises from three changes in perspective because of the non-circular and inclined orbit, the finite size of the Earth, and the orientation of the Moon in space. The first of these is called optical libration, the second parallax, and the third physical libration. Each of these can be divided into two contributions. The following are the three types of lunar libration: Optical libration, the combined libration of longitudinal and latitudinal libration produces a movement of the sub-Earth point and a wobbling view between the temporarily visible parts of the Moon, during a lunar orbit. This is not to be confused with the change of the Moon's apparent size because of the changing distance between the Moon and the Earth during the Moon's elliptic orbit, or with the change of positional angle because of the change in the position of the Moon's tilted axis, or with the observed swinging motion of the Moon because of the relative position of the Earth's tilted axis during an orbit of the Moon. Libration in longitude results from the eccentricity of the orbit of the Moon around the Earth; the Moon's rotation sometimes leads and sometimes lags its orbital position. The lunar libration in longitude was discovered by Johannes Hevelius in 1648. It can reach 7°54′ in amplitude. Longitudinal libration allows an observer on Earth to view at times further into the Moon's west and east respectively at different phases of the Moon's orbit. Libration in latitude results from the Moon's axial tilt (about 6.7°) between its rotation axis and orbital axis around Earth. This is analogous to how Earth's seasons arise from its axial tilt (about 23.4°) between its rotation axis and orbital axis about the Sun. Galileo Galilei is sometimes credited with the discovery of the lunar libration in latitude in 1632 although Thomas Harriot or William Gilbert might have done so before. Note Cassini's laws. It can reach 6°50′ in amplitude. The 6.7° depends on the orbit inclination of 5.15° and the negative equatorial tilt of 1.54°. Latitudinal libration allows an observer on Earth to view beyond the Moon's north pole and south pole at different phases of the Moon's orbit. Parallax libration depends on both the longitude and latitude of the location on Earth from which the Moon is observed. Diurnal libration is the small daily libration and oscillation from Earth's rotation, which carries an observer first to one side and then to the other side of the straight line joining Earth's and the Moon's centers, allowing the observer to look first around one side of the Moon and then around the other—since the observer is on Earth's surface, not at its center. It reaches less than 1° in amplitude. Physical libration is the oscillation of orientation in space about uniform rotation and precession. There are physical librations about all three axes. The sizes are roughly 100 seconds of arc. As seen from the Earth, this amounts to less than 1 second of arc. Forced physical librations can be predicted given the orbit and shape of the Moon. The periods of free physical librations can also be predicted, but their amplitudes and phases cannot be predicted. Physical libration Also called real libration, as opposed to the optical libration of longitudinal, latitudinal and diurnal types, the orientation of the Moon exhibits small oscillations of the pole direction in space and rotation about the pole. This libration can be differentiated between forced and free libration. Forced libration is caused by the forces exerted during the Moon's orbit around the Earth and the Sun, and free libration represents oscillations that occur over longer time periods. Forced physical libration Cassini's laws state the following: The Moon rotates uniformly about its polar axis keeping one side toward the Earth. The Moon's equator plane is tilted with respect to the ecliptic plane and it precesses uniformly along the ecliptic plane. The descending node of the equator on the ecliptic matches the ascending node of the orbit plane. In addition to uniform rotation and uniform precession of the equator plane, the Moon has small oscillations of orientation in space about all three axes. These oscillations are called physical librations. Apart from the 1.5427° tilt between equator and ecliptic, the oscillations are approximately ±100 seconds of arc in size. These oscillations can be expressed with trigonometric series that depend on the lunar moments of inertia A < B < C. The sensitive combinations are β = (C – A)/B and γ = (B – A)/C. The oscillation about the polar axis is most sensitive to γ and the 2-dimensional direction of the pole, including the 1.5427° tilt, is most sensitive to β. Consequently, accurate measurements of the physical librations provide accurate determinations of β = and γ = . The placement of three retroreflectors on the Moon by the Lunar Laser Ranging experiment and two retroreflectors by Lunokhod rovers allowed accurate measurement of the physical librations by laser ranging to the Moon. Free physical libration A free physical libration is similar to the solution of the reduced equation for linear differential equations. The periods of the free librations can be calculated, but their amplitudes must be measured. Lunar Laser Ranging provides the determinations. The two largest free librations were discovered by O. Calame. Modern values are: 1.3 seconds of arc with a 1056-day (2.9-year) period for rotation about the polar axis, a 74.6-year elliptical wobble of the pole of size 8.18 × 3.31 arcseconds, and an 81-year rotation of the pole in space that is 0.03 seconds of arc in size. The fluid core can cause a fourth mode with a period around four centuries. The free librations are expected to damp out in times very short compared to the age of the Moon. Consequently, their existence implies that there must be one or more stimulating mechanisms.	Physical sciences	Celestial mechanics	Astronomy
314510	https://en.wikipedia.org/wiki/Dire%20wolf	Dire wolf	The dire wolf (Aenocyon dirus ) is an extinct canine. The dire wolf lived in the Americas during the Late Pleistocene and Early Holocene epochs (125,000–9,500 years ago). A putative, controversial fossil was recently reported from northeast China, but other researchers questioned the taxonomic attribution of this specimen. The species was named in 1858, four years after the first specimen had been found. Two subspecies are recognized: Aenocyon dirus guildayi and Aenocyon dirus dirus. The largest collection of its fossils has been obtained from the Rancho La Brea Tar Pits in Los Angeles. Dire wolf remains have been found across a broad range of habitats including the plains, grasslands, and some forested mountain areas of North America and the arid savanna of South America. The sites range in elevation from sea level to . Dire wolf fossils have rarely been found north of 42°N latitude; there have been only five unconfirmed reports above this latitude. This range restriction is thought to be due to temperature, prey, or habitat limitations imposed by proximity to the Laurentide and Cordilleran ice sheets that existed at the time. The dire wolf was about the same size as the largest modern gray wolves (Canis lupus): the Yukon wolf and the northwestern wolf. A.d.guildayi weighed on average and A.d.dirus was on average . Its skull and dentition matched those of C.lupus, but its teeth were larger with greater shearing ability, and its bite force at the canine tooth was stronger than any known Canis species. These characteristics are thought to be adaptations for preying on Late Pleistocene megaherbivores, and in North America, its prey is known to have included western horses, ground sloths, mastodons, ancient bison, and camels. Its extinction occurred during the Quaternary extinction event along with its main prey species. Its reliance on megaherbivores has been proposed as the cause of its extinction, along with climatic change and competition with other species, or a combination of those factors. Dire wolves lived as recently as 9,500 years ago, according to dated remains. Taxonomy From the 1850s, the fossil remains of extinct large wolves were being found in the United States, and it was not immediately clear that these all belonged to one species. The first specimen of what would later become associated with Aenocyon dirus was found in mid-1854 in the bed of the Ohio River near Evansville, Indiana. The fossilized jawbone with cheek teeth was obtained by the geologist Joseph Granville Norwood from an Evansville collector, Francis A. Linck. The paleontologist Joseph Leidy determined that the specimen represented an extinct species of wolf and reported it under the name of Canis primaevus. Norwood's letters to Leidy are preserved along with the type specimen (the first of a species that has a written description) at the Academy of Natural Sciences of Philadelphia. In 1857, while exploring the Niobrara River valley in Nebraska, Leidy found the vertebrae of an extinct Canis species that he reported the following year under the name C.dirus. The name C.primaevus (Leidy 1854) was later renamed Canis indianensis (Leidy 1869) when Leidy found out that the name C.primaevus had previously been used by the British naturalist Brian Houghton Hodgson for the dhole. In 1876 the zoologist Joel Asaph Allen discovered the remains of Canis mississippiensis (Allen 1876) and associated these with C.dirus (Leidy 1858) and Canis indianensis (Leidy 1869). As so little was found of these three specimens, Allen thought it best to leave each specimen listed under its provisional name until more material could be found to reveal their relationship. In 1908 the paleontologist John Campbell Merriam began retrieving numerous fossilized bone fragments of a large wolf from the Rancho LaBrea tar pits. By 1912 he had found a skeleton sufficiently complete to be able to formally recognize these and the previously found specimens under the name C.dirus (Leidy 1858). Because the rules of nomenclature stipulated that the name of a species should be the oldest name ever applied to it, Merriam therefore selected the name of Leidy's 1858 specimen, C.dirus. In 1915 the paleontologist Edward Troxell indicated his agreement with Merriam when he declared C.indianensis a synonym of C.dirus. In 1918, after studying these fossils, Merriam proposed consolidating their names under the separate genus Aenocyon (from ainos, 'terrible' and cyon, 'dog') to become Aenocyon dirus, but at that time not everyone agreed with this extinct wolf being placed in a new genus separate from the genus Canis. Canis ayersi (Sellards 1916) and Aenocyon dirus (Merriam 1918) were recognized as synonyms of C.dirus by the paleontologist Ernest Lundelius in 1972. All of the above taxa were declared synonyms of C.dirus in 1979, according to the paleontologist Ronald M. Nowak. In 1984 a study by Björn Kurtén recognized a geographic variation within the dire wolf populations and proposed two subspecies: Canis dirus guildayi (named by Kurtén in honor of the paleontologist John E. Guilday) for specimens from California and Mexico that exhibited shorter limbs and longer teeth, and Canis dirus dirus for specimens east of the North American Continental Divide that exhibited longer limbs and shorter teeth. Kurtén designated a maxilla found in Hermit's Cave, New Mexico as representing the nominate subspecies C. d. dirus. In 2021, a DNA study found the dire wolf to be a highly divergent lineage when compared with the extant wolf-like canines, and this finding is consistent with the previously proposed taxonomic classification of the dire wolf as genus Aenocyon (Ancient Greek: "terrible wolf") as proposed by Merriam in 1918. Evolution In North America, the canid family came into existence 40 million years ago, and the canine subfamily Caninae about 32 million years ago. From the Caninae, the ancestors of the fox-like Vulpini and the dog-like Canini came into existence 9 million years ago. This group was first represented by Eucyon, and mostly by coyote-like Eucyon davisi that was spread widely across North America. From the Canini the Cerdocyonina, today represented by the South American canids, came into existence 6–5 million years ago. Its sister the wolf-like Canina came into existence 5 million years ago, however, they are likely to have originated as far back as 9 million years ago. Around 7 million years ago, the canines expanded into Eurasia and Africa, with Eucyon giving rise to the first of the genus Canis in Europe. Around 4–3 million years ago C. chihliensis, the first wolf-sized member of Canis, arose in China and expanded to give rise to other wolf-like members across Eurasia and Africa. Members of the genus Canis would later expand into North America. The dire wolf evolved in North America. However, its ancestral lineage is debated, with two competing theories. The first theory is based on fossil morphology, which indicates that an expansion of the genus Canis out of Eurasia led to the dire wolf. The second theory is based on DNA evidence, which indicates that the dire wolf arose from an ancestral lineage that originated in the Americas and was separate from the genus Canis. Morphological evidence Morphological evidence based on fossil remains indicates an expansion of genus Canis from out of Eurasia led to the dire wolf. In 1974 Robert A. Martin proposed that the large North American wolf C. armbrusteri (Armbruster's wolf) was C. lupus. Nowak, Kurtén, and Annalisa Berta proposed that C. dirus was not derived from C. lupus. In 1987, a new hypothesis proposed that a mammal population could give rise to a larger form called a hypermorph during times when food was abundant, but when food later became scarce the hypermorph would either adapt to a smaller form or go extinct. This hypothesis might explain the large body sizes found in many Late Pleistocene mammals compared to their modern counterparts. Both extinction and speciationa process by which a new species splits from an older onecould occur together during periods of climatic extremes. Gloria D. Goulet agreed with Martin, proposing further that this hypothesis might explain the sudden appearance of C. dirus in North America and, judging from the similarities in their skull shapes, that C. lupus had given rise to the C. dirus hypermorph due to an abundance of game, a stable environment, and large competitors. The three paleontologists Xiaoming Wang, Richard H. Tedford, and Ronald M. Nowak propose that C. dirus evolved from Canis armbrusteri, with Nowak stating that both species arose in the Americas and that specimens found in Cumberland Cave, Maryland, appear to be C. armbrusteri diverging into C. dirus. Nowak believed that Canis edwardii was the first appearance of the wolf in North America, and it appears to be close to the lineage which produced C. armbrusteri and C. dirus. Tedford believes that the early wolf from China, Canis chihliensis, may have been the ancestor of both C. armbrusteri and the gray wolf C. lupus. The sudden appearance of C. armbrusteri in mid-latitude North America during the Early Pleistocene 1.5 million years ago, along with the mammoth, suggests that it was an immigrant from Asia, with the gray wolf C. lupus evolving in Beringia later in the Pleistocene and entering mid-latitude North America during the Last Glacial Period along with its Beringian prey. In 2010 Francisco Prevosti proposed that C. dirus was a sister taxon to C. lupus. C. dirus lived in the Late Pleistocene to the early Holocene, 125,000–10,000 YBP (years before present), in North and South America. The majority of fossils from the eastern C. d. dirus have been dated 125,000–75,000 YBP, but the western fossils are not only smaller in size but more recent; thus it has been proposed that derived from However, there are disputed specimens of C. dirus that date to 250,000 YBP. Fossil specimens of C. dirus discovered at four sites in the Hay Springs area of Sheridan County, Nebraska, were named Aenocyon dirus nebrascensis (Frick 1930, undescribed), but Frick did not publish a description of them. Nowak later referred to this material as C. armbrusteri; then, in 2009, Tedford formally published a description of the specimens and noted that, although they exhibited some morphological characteristics of both C. armbrusteri and C. dirus, he referred to them only as C. dirus. A fossil discovered in the Horse Room of the Salamander Cave in the Black Hills of South Dakota may possibly be C. dirus; if so, this fossil is one of the earliest specimens on record. It was catalogued as Canis cf. C. dirus (where cf. in Latin means confer, uncertain). The fossil of a horse found in the Horse Room provided a uranium-series dating of 252,000 YBP and the Canis cf. dirus specimen was assumed to be from the same period. C. armbrusteri and C. dirus share some characteristics (synapomorphies) that imply the latter's descent from the former. The fossil record suggests C. dirus originated around 250,000 YBP in the open terrain of the mid-continent before expanding eastward and displacing its ancestor C. armbrusteri. The first appearance of C. dirus would therefore be 250,000 YBP in California and Nebraska, and later in the rest of the United States, Canada, Mexico, Venezuela, Ecuador, Bolivia, and Peru, but the identity of these earliest fossils is not confirmed. In South America, C. dirus specimens dated to the Late Pleistocene were found along the north and west coasts, but none have been found in Argentina, an area that was inhabited by Canis gezi and Canis nehringi. Given their similarities and timeframes, it is proposed that C. gezi was the ancestor of Canis nehringi. One study found that C. dirus was more evolutionarily derived compared with C. nehringi, and was larger in the size and construction of its lower molars for more efficient predation. For this reason, some researchers have proposed that C. dirus may have originated in South America. Tedford proposed that C. armbrusteri was the common ancestor for both the North and South American wolves. Later studies concluded that C. dirus and C. nehringi were the same species, and that C. dirus had migrated from North America into South America, making it a participant in the Great American Interchange. In 2018, a study found that Canis gezi did not fall under genus Canis and should be classified under the subtribe Cerdocyonina, however no genus was proposed. The 2020 discovery of a claimed dire wolf fossil in northeast China indicates that dire wolves may have crossed Beringia when it existed, though other researchers doubt that this specimen represents a dire wolf. DNA evidence DNA evidence indicates the dire wolf arose from an ancestral lineage that originated in the Americas and was separate to genus Canis. In 1992 an attempt was made to extract a mitochondrial DNA sequence from the skeletal remains of A.d.guildayi to compare its relationship to other Canis species. The attempt was unsuccessful because these remains had been removed from the LaBrea pits and tar could not be removed from the bone material. In 2014 an attempt to extract DNA from a Columbian mammoth from the tar pits also failed, with the study concluding that organic compounds from the asphalt permeate the bones of all ancient samples from the LaBrea pits, hindering the extraction of DNA samples. In 2021, researchers sequenced the nuclear DNA (from the cell nucleus) taken from five dire wolf fossils dating from 13,000 to 50,000 years ago. The sequences indicate the dire wolf to be a highly divergent lineage which last shared a most recent common ancestor with the wolf-like canines 5.7 million years ago. The study also measured numerous dire wolf and gray wolf skeletal samples that showed their morphologies to be highly similar, which had led to the theory that the dire wolf and the gray wolf had a close evolutionary relationship. The morphological similarity between dire wolves and gray wolves was concluded to be due to convergent evolution. Members of the wolf-like canines are known to hybridize with each other but the study could find no indication of genetic admixture from the five dire wolf samples with extant North American gray wolves and coyotes nor their common ancestor. This finding indicates that the wolf and coyote lineages evolved in isolation from the dire wolf lineage. The study proposes an early origin of the dire wolf lineage in the Americas, and that this geographic isolation allowed them to develop a degree of reproductive isolation since their divergence 5.7 million years ago. Coyotes, dholes, gray wolves, and the extinct Xenocyon evolved in Eurasia and expanded into North America relatively recently during the Late Pleistocene, therefore there was no admixture with the dire wolf. The long-term isolation of the dire wolf lineage implies that other American fossil taxa, including C. armbrusteri and C. edwardii, may also belong to the dire wolf's lineage. The study's findings are consistent with the previously proposed taxonomic classification of the dire wolf as genus Aenocyon. Radiocarbon dating The age of most dire wolf localities is determined solely by biostratigraphy, but biostratigraphy is an unreliable indicator within asphalt deposits. Some sites have been radiocarbon dated, with dire wolf specimens from the LaBrea pits dated in calendar years as follows: 82 specimens dated 13,000–14,000YBP; 40 specimens dated 14,000–16,000YBP; 77 specimens dated 14,000–18,000YBP; 37 specimens dated 17,000–18,000YBP; 26 specimens dated 21,000–30,000YBP; 40 specimens dated 25,000–28,000YBP; and 6specimens dated 32,000–37,000YBP. A specimen from Powder Mill Creek Cave, Missouri, was dated at 13,170YBP. Description The average dire wolf proportions were similar to those of two modern North American wolves: the Yukon wolf (Canis lupus pambasileus) and the Northwestern wolf (Canis lupus occidentalis). The largest northern wolves today have a shoulder height of up to and a body length of . Some dire wolf specimens from Rancho LaBrea are smaller than this, and some are larger. The dire wolf had smaller feet and a larger head when compared with a northern wolf of the same body size. The skull length could reach up to or longer, with a broader palate, frontal region, and zygomatic arches compared with the Yukon wolf. These dimensions make the skull very massive. Its sagittal crest was higher, with the inion showing a significant backward projection, and with the rear ends of the nasal bones extending relatively far back into the skull. A connected skeleton of a dire wolf from Rancho LaBrea is difficult to find because the tar allows the bones to disassemble in many directions. Parts of a vertebral column have been assembled, and it was found to be similar to that of the modern wolf, with the same number of vertebrae. Geographic differences in dire wolves were not detected until 1984, when a study of skeletal remains showed differences in a few cranio-dental features and limb proportions between specimens from California and Mexico (A.d.guildayi) and those found from the east of the Continental Divide (A.d.dirus). A comparison of limb size shows that the rear limbs of A.d.guildayi were 8% shorter than the Yukon wolf due to a significantly shorter tibia and metatarsus, and that the front limbs were also shorter due to their slightly shorter lower bones. With its comparatively lighter and smaller limbs and massive head, A.d.guildayi was not as well adapted for running as timber wolves and coyotes. A.d.dirus possessed significantly longer limbs than A.d.guildayi. The forelimbs were 14% longer than A.d.guildayi due to 10% longer humeri, 15% longer radii, and 15% longer metacarpals. The rear limbs were 10% longer than A.d.guildayi due to 10% longer femora and tibiae, and 15% longer metatarsals. A.d.dirus is comparable to the Yukon wolf in limb length. The largest A.d.dirus femur was found in Carroll Cave, Missouri, and measured . A.d.guildayi is estimated to have weighed on average , and A.d.dirus weighed on average with some specimens being larger, but these could not have exceeded due to skeletal limits. In comparison, the average weight of the Yukon wolf is for males and for females. Individual weights for Yukon wolves can vary from to , with one Yukon wolf weighing . These figures show the average dire wolf to be similar in size to the largest modern gray wolf. The remains of a complete male A. dirus are sometimes easy to identify compared to other Canis specimens because the baculum (penis bone) of the dire wolf is very different from that of all other living canids. A 2024 study found the baculum of a male dire wolf to be proportionally longer than the baculum of modern canids, which may be indicative of stronger competition between males and unusual behaviors among canids including non-monogamous mating. Adaptation Ecological factors such as habitat type, climate, prey specialization, and predatory competition have been shown to greatly influence gray wolf craniodental plasticity, which is an adaptation of the cranium and teeth due to the influences of the environment. Similarly, the dire wolf was a hypercarnivore, with a skull and dentition adapted for hunting large and struggling prey; the shape of its skull and snout changed across time, and changes in the size of its body have been correlated with climate fluctuations. Paleoecology The last glacial period, commonly referred to as the "Ice Age", spanned 125,000–14,500YBP and was the most recent glacial period within the current ice age, which occurred during the last years of the Pleistocene era. The Ice Age reached its peak during the Last Glacial Maximum, when ice sheets began advancing from 33,000YBP and reached their maximum limits 26,500YBP. Deglaciation commenced in the Northern Hemisphere approximately 19,000YBP and in Antarctica approximately 14,500YBP, which is consistent with evidence that glacial meltwater was the primary source for an abrupt rise in sea level 14,500YBP. Access into northern North America was blocked by the Wisconsin glaciation. The fossil evidence from the Americas points to the extinction mainly of large animals, termed Pleistocene megafauna, near the end of the last glaciation. Coastal southern California from 60,000YBP to the end of the Last Glacial Maximum was cooler and with a more balanced supply of moisture than today. During the Last Glacial Maximum, the mean annual temperature decreased from down to degrees, and annual precipitation had decreased from down to . This region was unaffected by the climatic effects of the Wisconsin glaciation and is thought to have been an Ice Age refugium for animals and cold-sensitive plants. By 24,000YBP, the abundance of oak and chaparral decreased, but pines increased, creating open parklands similar to today's coastal montane/juniper woodlands. After 14,000YBP, the abundance of conifers decreased, and those of the modern coastal plant communities, including oak woodland, chaparral, and coastal sage scrub, increased. The Santa Monica Plain lies north of the city of Santa Monica and extends along the southern base of the Santa Monica Mountains, and 28,000–26,000YBP it was dominated by coastal sage scrub, with cypress and pines at higher elevations. The Santa Monica Mountains supported a chaparral community on its slopes and isolated coast redwood and dogwood in its protected canyons, along with river communities that included willow, red cedar, and sycamore. These plant communities suggest a winter rainfall similar to that of modern coastal southern California, but the presence of coast redwood now found to the north indicates a cooler, moister, and less seasonal climate than today. This environment supported large herbivores that were prey for dire wolves and their competitors. Prey A range of animal and plant specimens that became entrapped and were then preserved in tar pits have been removed and studied so that researchers can learn about the past. The Rancho LaBrea tar pits located near Los Angeles in Southern California are a collection of pits of sticky asphalt deposits that differ in deposition time from 40,000 to 12,000YBP. Commencing 40,000YBP, trapped asphalt has been moved through fissures to the surface by methane pressure, forming seeps that can cover several square meters and be deep. A large number of dire wolf fossils have been recovered from the La Brea tar pits. Over 200,000 specimens (mostly fragments) have been recovered from the tar pits, with the remains ranging from Smilodon to squirrels, invertebrates, and plants. The time period represented in the pits includes the Last Glacial Maximum when global temperatures were lower than today, the Pleistocene–Holocene transition (Bølling-Allerød interval), the Oldest Dryas cooling, the Younger Dryas cooling from 12,800 to 11,500YBP, and the American megafaunal extinction event 12,700YBP when 90 genera of mammals weighing over became extinct. Isotope analysis can be used to identify some chemical elements, allowing researchers to make inferences about the diet of the species found in the pits. Isotope analysis of bone collagen extracted from LaBrea specimens provides evidence that the dire wolf, Smilodon, and the American lion (Panthera atrox) competed for the same prey. Their prey included the extinct camel Camelops hesternus, the extinct bison Bison antiquus, the "dwarf" pronghorn (Capromeryx minor), the equine Equus occidentalis, and Harlan's ground sloth (Paramylodon harlani) native to North American grasslands. The Columbian mammoth (Mammuthus columbi) and the American mastodon (Mammut americanum) were rare at LaBrea. The horses remained mixed feeders and the pronghorns mixed browsers, but at the Last Glacial Maximum and its associated shift in vegetation the camels and bison were forced to rely more heavily on conifers. A study of isotope data of La Brea dire wolf fossils dated 10,000YBP provides evidence that the horse was an important prey species at the time, and that sloth, mastodon, bison, and camel were less common in the dire wolf diet. This indicates that the dire wolf was not a prey specialist, and at the close of the Late Pleistocene before its extinction it was hunting or scavenging the most available herbivores. A study based on specimens found in Cedral, San Luis Potosi found that the dire wolf primarily preyed on herbivores that consumed C4 plants and on mixed-diet herbivores. Dentition and bite force When compared with the dentition of genus Canis members, the dire wolf was considered the most evolutionary derived (advanced) wolf-like species in the Americas. The dire wolf could be identified separately from all other Canis species by its possession of: "P2 with a posterior cusplet; P3 with two posterior cusplets; M1 with a mestascylid, entocristed, entoconulid, and a transverse crest extending from the metaconid to the hyperconular shelf; M2 with entocristed and entoconulid." A study of the estimated bite force at the canine teeth of a large sample of living and fossil mammalian predators, when adjusted for the body mass, found that for placental mammals the bite force at the canines (in newtons/kilogram of body weight) was greatest in the dire wolf (163), followed among the modern canids by the four hypercarnivores that often prey on animals larger than themselves: the African hunting dog (142), the gray wolf (136), the dhole (112), and the dingo (108). The bite force at the carnassials showed a similar trend to the canines. A predator's largest prey size is strongly influenced by its biomechanical limits. The morphology of the dire wolf was similar to that of its living relatives, and assuming that the dire wolf was a social hunter, then its high bite force relative to living canids suggests that it preyed on relatively large animals. The bite force rating of the bone-consuming spotted hyena (117) challenged the common assumption that high bite force in the canines and the carnassials was necessary to consume bone. A study of the cranial measurements and jaw muscles of dire wolves found no significant differences with modern gray wolves in all but 4 of 15 measures. Upper dentition was the same except that the dire wolf had larger dimensions, and the P4 had a relatively larger, more massive blade that enhanced slicing ability at the carnassial. The jaw of the dire wolf had a relatively broader and more massive temporalis muscle, able to generate slightly more bite force than the gray wolf. Due to the jaw arrangement, the dire wolf had less temporalis leverage than the gray wolf at the lower carnassial (m1) and lower p4, but the functional significance of this is not known. The lower premolars were relatively slightly larger than those of the gray wolf, and the dire wolf m1 was much larger and had more shearing ability. The dire wolf canines had greater bending strength than those of living canids of equivalent size and were similar to those of hyenas and felids. All these differences indicate that the dire wolf was able to deliver stronger bites than the gray wolf, and with its flexible and more rounded canines was better adapted for struggling with its prey. Behavior At La Brea, predatory birds and mammals were attracted to dead or dying herbivores that had become mired, and then these predators became trapped themselves. Herbivore entrapment was estimated to have occurred once every fifty years, and for every instance of herbivore remains found in the pits there were an estimated ten carnivores. A.d.guildayi is the most common carnivoran found at LaBrea, followed by Smilodon. Remains of dire wolves outnumber remains of gray wolves in the tar pits by a ratio of five to one. During the Last Glacial Maximum, coastal California, with a climate slightly cooler and wetter than today, is thought to have been a refuge, and a comparison of the frequency of dire wolves and other predator remains at LaBrea to other parts of California and North America indicates significantly greater abundances; therefore, the higher dire wolf numbers in the LaBrea region did not reflect the wider area. Assuming that only a few of the carnivores that were feeding became trapped, it is likely that fairly sizeable groups of dire wolves fed together on these occasions. The difference between the male and female of a species apart from their sex organs is called sexual dimorphism, and in this regard little variance exists among the canids. A study of dire wolf remains dated 15,360–14,310YBP and taken from one pit that focused on skull length, canine tooth size, and lower molar length showed little dimorphism, similar to that of the gray wolf, indicating that dire wolves lived in monogamous pairs. Their large size and highly carnivorous dentition supports the proposal that the dire wolf was a predator that fed on large prey. To kill ungulates larger than themselves, the African wild dog, the dhole, and the gray wolf depend on their jaws as they cannot use their forelimbs to grapple with prey, and they work together as a pack consisting of an alpha pair and their offspring from the current and previous years. It can be assumed that dire wolves lived in packs of relatives that were led by an alpha pair. Large and social carnivores would have been successful at defending carcasses of prey trapped in the tar pits from smaller solitary predators, and thus the most likely to become trapped themselves. The many A.d.guildayi and Smilodon remains found in the tar pits suggests that both were social predators. All social terrestrial mammalian predators prey mostly on terrestrial herbivorous mammals with a body mass similar to the combined mass of the social group members attacking the prey animal. The large size of the dire wolf provides an estimated prey size in the range. Stable isotope analysis of dire wolf bones provides evidence that they had a preference for consuming ruminants such as bison rather than other herbivores but moved to other prey when food became scarce, and occasionally scavenged on beached whales along the Pacific coast when available. A pack of timber wolves can bring down a moose that is their preferred prey, and a pack of dire wolves bringing down a bison is conceivable. Although some studies have suggested that because of tooth breakage, the dire wolf must have gnawed bones and may have been a scavenger, its widespread occurrence and the more gracile limbs of the dire wolf indicate a predator. Like the gray wolf today, the dire wolf probably used its post-carnassial molars to gain access to marrow, but the dire wolf's larger size enabled it to crack larger bones. Tooth breakage Tooth breakage is related to a carnivore's behavior. A study of nine modern carnivores found that one in four adults had suffered tooth breakage and that half of these breakages were of the canine teeth. The most breakage occurred in the spotted hyena that consumes all of its prey including the bone; the least breakage occurred in the African wild dog, and the gray wolf ranked in between these two. The eating of bone increases the risk of accidental fracture due to the relatively high, unpredictable stresses that it creates. The most commonly broken teeth are the canines, followed by the premolars, carnassial molars, and incisors. Canines are the teeth most likely to break because of their shape and function, which subjects them to bending stresses that are unpredictable in both direction and magnitude. The risk of tooth fracture is also higher when killing large prey. A study of the fossil remains of large carnivores from LaBrea pits dated 36,000–10,000YBP shows tooth breakage rates of 5–17% for the dire wolf, coyote, American lion, and Smilodon, compared to 0.5–2.7% for ten modern predators. These higher fracture rates were across all teeth, but the fracture rates for the canine teeth were the same as in modern carnivores. The dire wolf broke its incisors more often when compared to the modern gray wolf; thus, it has been proposed that the dire wolf used its incisors more closely to the bone when feeding. Dire wolf fossils from Mexico and Peru show a similar pattern of breakage. A 1993 study proposed that the higher frequency of tooth breakage among Pleistocene carnivores compared with living carnivores was not the result of hunting larger game, something that might be assumed from the larger size of the former. When there is low prey availability, the competition between carnivores increases, causing them to eat faster and thus consume more bone, leading to tooth breakage. As their prey became extinct around 10,000 years ago, so did these Pleistocene carnivores, except for the coyote (which is an omnivore). A later La Brea pits study compared tooth breakage of dire wolves in two time periods. One pit contained fossil dire wolves dated 15,000YBP and another dated 13,000YBP. The results showed that the 15,000YBP dire wolves had three times more tooth breakage than the 13,000YBP dire wolves, whose breakage matched those of nine modern carnivores. The study concluded that between 15,000 and 14,000YBP prey availability was less or competition was higher for dire wolves and that by 13,000YBP, as the prey species moved towards extinction, predator competition had declined and therefore the frequency of tooth breakage in dire wolves had also declined. Carnivores include both pack hunters and solitary hunters. The solitary hunter depends on a powerful bite at the canine teeth to subdue their prey, and thus exhibits a strong mandibular symphysis. In contrast, a pack hunter, which delivers many shallower bites, has a comparably weaker mandibular symphysis. Thus, researchers can use the strength of the mandibular symphysis in fossil carnivore specimens to determine what kind of hunter it wasa pack hunter or a solitary hunterand even how it consumed its prey. The mandibles of canids are buttressed behind the carnassial teeth to enable the animals to crack bones with their post-carnassial teeth (molars M2 and M3). A study found that the mandible buttress profile of the dire wolf was lower than that of the gray wolf and the red wolf, but very similar to the coyote and the African hunting dog. The dorsoventrally weak symphyseal region (in comparison to premolars P3 and P4) of the dire wolf indicates that it delivered shallow bites similar to its modern relatives and was therefore a pack hunter. This suggests that the dire wolf may have processed bone but was not as well adapted for it as was the gray wolf. The fact that the incidence of fracture for the dire wolf reduced in frequency in the Late Pleistocene to that of its modern relatives suggests that reduced competition had allowed the dire wolf to return to a feeding behavior involving a lower amount of bone consumption, a behavior for which it was best suited. The results of a study of dental microwear on tooth enamel for specimens of the carnivore species from LaBrea pits, including dire wolves, suggest that these carnivores were not food-stressed just before their extinction. The evidence also indicated that the extent of carcass utilization (i.e., amount consumed relative to the maximum amount possible to consume, including breakup and consumption of bones) was less than among large carnivores today. These findings indicates that tooth breakage was related to hunting behavior and the size of prey. Climate impact Past studies proposed that changes in dire wolf body size correlated with climate fluctuations. A later study compared dire wolf craniodental morphology from four LaBrea pits, each representing four different time periods. The results are evidence of a change in dire wolf size, dental wear and breakage, skull shape, and snout shape across time. Dire wolf body size had decreased between the start of the Last Glacial Maximum and near its ending at the warm Allerød oscillation. Evidence of food stress (food scarcity leading to lower nutrient intake) is seen in smaller body size, skulls with a larger cranial base, and shorter snout (shape neoteny and size neoteny), and more tooth breakage and wear. Dire wolves dated 17,900YBP showed all of these features, which indicates food stress. Dire wolves dated 28,000YBP also showed to a degree many of these features but were the largest wolves studied, and it was proposed that these wolves were also suffering from food stress and that wolves earlier than this date were even bigger in size. Nutrient stress is likely to lead to stronger bite forces to more fully consume carcasses and to crack bones, and with changes to skull shape to improve mechanical advantage. North American climate records reveal cyclic fluctuations during the glacial period that included rapid warming followed by gradual cooling, called Dansgaard–Oeschger events. These cycles would have caused increased temperature and aridity, and at LaBrea would have caused ecological stress and therefore food stress. A similar trend was found with the gray wolf, which in the Santa Barbara basin was originally massive, robust, and possibly convergent evolution with the dire wolf, but was replaced by more gracile forms by the start of the Holocene. {\| class="wikitable" style="text-align:center;" \|+ Dire wolf information based on skull measurements !scope="col"\| Variable !scope="col"\| 28,000 YBP !scope="col"\| 26,100 YBP !scope="col"\| 17,900 YBP !scope="col"\| 13,800 YBP \|- !scope="row"\| Body size \| largest \| large \| smallest \| medium/small \|- !scope="row"\| Tooth breakage \| high \| low \| high \| low \|- !scope="row"\| Tooth wear \| high \| low \| high \| low \|- !scope="row"\| Snout shape \| shortening, largest cranial base \| average \| shortest, largest cranial base \| average \|- !scope="row"\| Tooth row shape \| robust \| – \| – \| gracile \|- !scope="row"\| DO event \| number 3 or 4 \| none \| imprecise data \| imprecise data \|} Competitors Just before the appearance of the dire wolf, North America was invaded by the Canis subgenus Xenocyon (ancestor of the Asian dhole and the African hunting dog) that was as large as the dire wolf and more hypercarnivorous. The fossil record shows them as rare, and it is assumed that they could not compete with the newly derived dire wolf. Stable isotope analysis provides evidence that the dire wolf, Smilodon, and the American lion competed for the same prey. Other large carnivores included the extinct North American giant short-faced bear (Arctodus simus), the modern cougar (Puma concolor), the Pleistocene coyote (Canis latrans), and the Pleistocene gray wolf that was more massive and robust than today. These predators may have competed with humans who hunted for similar prey. Specimens that have been identified by morphology as Beringian wolves (C.lupus) and radiocarbon dated 25,800–14,300 YBP have been found in the Natural Trap Cave at the base of the Bighorn Mountains in Wyoming, in the western United States. The location is directly south of what would at that time have been a division between the Laurentide Ice Sheet and the Cordilleran Ice Sheet. A temporary channel between the glaciers may have existed that allowed these large, Alaskan direct competitors of the dire wolf, which were also adapted for preying on megafauna, to come south of the ice sheets. Dire wolf remains are absent north of the 42°Nlatitude in North America, therefore, this region would have been available for Beringian wolves to expand south along the glacier line. How widely they were then distributed is not known. These also became extinct at the end of the Late Pleistocene, as did the dire wolf. After arriving in eastern Eurasia, the dire wolf would have likely faced competition from the area's most dominant, widespread predator, the eastern subspecies of cave hyena (Crocuta crocuta ultima). Competition with this species may have kept Eurasian dire wolf populations very low, leading to the paucity of dire wolf fossil remains in this otherwise well-studied fossil fauna. Range Dire wolf remains have been found across a broad range of habitats including the plains, grasslands, and some forested mountain areas of North America, the arid savannah of South America, and possibly the steppes of eastern Asia. The sites range in elevation from sea level to . The location of these fossil remains suggests that dire wolves lived predominantly in the open lowlands along with their prey the large herbivores. Dire wolf remains are not often found at high latitudes in North America, with the northernmost record in southern Canada. In the United States, dire wolf fossils have been reported in Arizona, California, Florida, Idaho, Indiana, Kansas, Kentucky, Missouri, Nebraska, New Mexico, Oregon, Pennsylvania, South Carolina, South Dakota, Texas, Utah, Virginia, West Virginia, Wyoming, and Nevada. The identity of fossils reported farther north than California is not confirmed. There have been five reports of unconfirmed dire wolf fossils north of 42°Nlatitude at Fossil Lake, Oregon (125,000–10,000YBP), American Falls Reservoir, Idaho (125,000–75,000YBP), Salamander Cave, South Dakota (250,000YBP), and four closely grouped sites in northern Nebraska (250,000YBP). This suggests a range restriction on dire wolves due to temperature, prey, or habitat. The major fossil-producing sites for A.d.dirus are located east of the Rocky Mountains and include Friesenhahn Cave, near San Antonio, Texas; Carroll Cave, near Richland, Missouri; and Reddick, Florida. Localities in Mexico where dire wolf remains have been collected include ElCedazo in Aguascalientes, Comondú Municipality in Baja California Sur, ElCedral in San Luis Potosí, ElTajo Quarry near Tequixquiac, state of Mexico, Valsequillo in Puebla, Lago de Chapala in Jalisco, Loltun Cave in Yucatán, Potrecito in Sinaloa, San Josecito Cave near Aramberri in Nuevo León and Térapa in Sonora. The specimens from Térapa were confirmed as A.d.guildayi. The finds at San Josecito Cave and ElCedazo have the greatest number of individuals from a single locality. In South America, dire wolves have been dated younger than 17,000 YBP and have been reported from six localities: Muaco in the western Falcón state of Venezuela, Talara Province in Peru, Monagas state in eastern Venezuela, the Tarija Department in Bolivia, Atacama Desert of Chile, and Ecuador. If the dire wolf originated in North America, the species likely dispersed into South America via the Andean corridor, a proposed pathway for temperate mammals to migrate from Central to South America because of the favorable cool, dry, and open habitats that characterized the region at times. This most likely happened during a glacial period because the pathway then consisted of open, arid regions and savanna, whereas during inter-glacial periods it would have consisted of tropical rain forest. In 2020, a fossil mandible (IVPP V25381) later analyzed as a dire wolf's was found in the vicinity of Harbin, northeastern China. The fossil was taxonomically described and dated 40,000 YBP. This discovery challenges previous theories that the cold temperatures and ice sheets at northern latitudes in North America would be a barrier for dire wolves, which was based on no dire wolf fossils being found above the 42° latitude in North America. It is proposed that the dire wolf followed migrating prey from mid-latitude North America then across Beringia into Eurasia. However, the 2022 study argued that the morphology and size of the specimen is inconclusive for its taxonomic determination as a dire wolf. Extinction During the Quaternary extinction event around 12,700YBP, 90genera of mammals weighing over became extinct. The extinction of the large carnivores and scavengers is thought to have been caused by the extinction of the megaherbivore prey upon which they depended. The cause of the extinction of the megafauna is debated but has been attributed to the impact of climatic change, competition with other species including overexploitation by newly arrived human hunters, or a combination of both. One study proposes that several extinction models should be investigated because so little is known about the biogeography of the dire wolf and its potential competitors and prey, nor how all these species interacted and responded to the environmental changes that occurred at the time of extinction. Ancient DNA and radiocarbon data indicate that local genetic populations were replaced by others within the same species or by others within the same genus. Both the dire wolf and the Beringian wolf went extinct in North America, leaving only the less carnivorous and more gracile form of the wolf to thrive, which may have outcompeted the dire wolf. One study proposes an early origin of the dire wolf lineage in the Americas which led to its reproductive isolation, such that when coyotes, dholes, gray wolves, and Xenocyon expanded into North America from Eurasia in the Late Pleistocene there could be no admixture with the dire wolf. Gray wolves and coyotes may have survived due to their ability to hybridize with other canids – such as the domestic dog – to acquire traits that resist diseases brought by taxa arriving from Eurasia. Reproductive isolation may have prevented the dire wolf from acquiring these traits. A 2023 study documented a high degree of subchondral defects in joint surfaces of dire wolf and Smilodon specimens from the La Brea Tar pits that resembled osteochondrosis dissecans. As modern dogs with this disease are inbred, the researchers suggested this would have been the case for the prehistoric species as well as they approached extinction, but cautioned that more research was needed to determine if this was also the case in specimens from other parts of the Americas. Dire wolf remains having the youngest geological ages are dated at 9,440YBP at Brynjulfson Cave, Boone County, Missouri, 9,860YBP at Rancho La Brea, California, and 10,690YBP at La Mirada, California. Dire wolf remains have been radiocarbon dated to 8,200YBP from Whitewater Draw in Arizona, though one author has stated that radiocarbon dating of bone carbonate is unreliable. All of these dates are uncalibrated and the actual age of the remains is likely older. In South America, the most recent remains at Talara, Peru date to 9,030 ± 240 YBP (also uncalibrated), while the most recent remains of "C. nehringi" from Luján, Argentina are older than the most recent stratigraphical section of the site, dated to 10–11,000 YBP.	Biology and health sciences	Canines	Animals

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